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MALDI-TOF Mass Spectrometry

JoVE 5691

Matrix-assisted laser desorption ionization (MALDI) is a mass spectrometry ion source ideal for the analysis of biomolecules. Instead of ionizing compounds in the gaseous state, samples are embedded in a matrix, which is struck by a laser. The matrix absorbs the majority of the energy; some of this energy is then transferred to the sample, which ionizes as a result. Sample ions can then be identified using a time-of-flight analyzer (TOF).

This video covers principles of MALDI-TOF, including matrix selection and how TOF is used to elucidate mass-to-charge ratios. This procedure shows the preparation of a MALDI plate, the loading of samples onto the plate, and the operation of the TOF-mass spectrometer. In the final section, applications and variations are shown, including whole-cell analysis, characterization of complex biological samples, and electron spray ionization.

Matrix-assisted laser desorption ionization, or MALDI, is a mass spectrometry ion source ideal for the analysis of biomolecules. Most ion sources remove structural information from large, fragile biomolecules. MALDI maintains structural integrity, and therefore information, while accelerating the molecules into the mass analyzer, which separates the compounds based on size and charge. The most commonly coupled with MALDI is the time of flight, or TOF, mass analyzer. This video will show the concepts of MALDI ionization, a general procedure, and some of its uses in biochemistry.

For mass spectrometry to function, molecules must be ionized into the gaseous state. In MALDI, the sample is embedded in a matrix, typically an organic compound containing aromatic and conjugated double bonds.

When a laser pulse strikes this mixture the matrix absorbs the majority of the energy, rapidly heats, and is desorbed, or released, from the surface. The energized matrix transfers some of its energy to the biomolecules, desorbing and then ionizing them.

MALDI is typically paired with a time of flight, or TOF, mass analyzer. An electric field applies kinetic energy to the ions, moving them into a field-free region called a drift tube. The velocity of the ions as they move through the tube is related to their mass-to-charge ratio, so heavier particles travel slower through instrument. A detector at the end of the tube measures each ion's flight time. With this knowledge, as well as the tube length and applied field strength, the mass-to-charge ratio of each ion can be elucidated.

This plot of signal intensity to mass-to-charge-ratio, known as a mass spectrum, can be compared to a library of collected spectra. If no matches are found, it can molecules can be identified by further techniques, such as tandem mass spectrometry. For more information, see this collection's video on the topic.

Now that the basics of MALDI-TOF have been discussed, let's look at the process in the laboratory.

Before beginning an experiment, it's important to consider the choice of matrix from which samples will be desorbed. It must absorb the laser energy, be stable in a vacuum, not react with the target molecules, and be able to desorb. Depending on the sample, different matrices are preferred. For a large protein, a combination of CHCA and DHB has shown better separation of the peaks, called resolution, than the individual matrices.

There are a number of ways to prepare samples. We'll show what is known as the "double-layer", or "sandwich," method. To begin, clean the MALDI plate with ultra-pure reagents, as mass spectrometry is very sensitive to contamination. Dry the plate with a stream of inert gas.

Next, a saturated matrix solution is made, typically with an organic solvent . The solution is streaked onto the MALDI plate and dried. A second saturated solution of matrix containing trifluoroacetic acid, or TFA, is prepared. TFA helps ions into the gaseous phase.

Next, the sample solution is added on top of the dried matrix spot. Add the matrix solution containing TFA on top of the sample, thereby completing the matrix "sandwich". Homogeneity of the spot can be verified under a low-powered microscope.

Plate a calibration standard, which is a mixture with a wide range of known masses and is used to correlate the time-of-flight to m/z. Finally, plate the matrix alone as a negative control.

To analyze the spots, place the target plate into the instrument. Ensure there's no debris present, allowing for the formation of a tight vacuum. In the software, select the standard, negative control, and samples of interest. Label the spots with the correct identification.

The ion source and lens voltages can be manipulated to improve performance of the analysis. This will depend on the specifics of the instrument and sample. Focus on the standard spot and calibrate the instrument with the software.

Next, collect spectra from each of the sample spots. Try a few different locations on the spot to maximize the quality of the collected data. Once finished, the MALDI plate can be collected and reused after cleaning.

Now that we've reviewed a procedure, let's look at some of the ways MALDI is utilized, and a different ionization technique.

In addition to biomolecules, MALDI can be used to analyze living cells. Macrophages are immune cells that take on one of several different forms, based on their microenvironment. After exposing the cells to various signaling molecules, or cytokines, they can be added directly to the plate, and analyzed. The MALDI spectra can be used as unique "fingerprints", depending on the cytokine used.

Complex biological samples like mammalian sebaceous secretions require a step of purification before MALDI analysis. Thin layer chromatography is one such technique that relies on the components' polarity. The compounds are collected from the TLC pate, purified, and transferred to a MALDI matrix. The resulting spectra verify the identity and purity of the separated biomolecules from the mammalian sebaceous secretions.

Another common ion source for biomolecules is electrospray ionization, or ESI. In this method, the sample is injected into the instrument, where a high voltage is applied, creating an aerosol of charged droplets. As the solvent in the droplet evaporates, the charge is moved to the sample molecules, till they are completely gaseous. ESI doesn't require the spotting procedure, and the sample can be injected directly into the instrument. On the other hand, ESI is more sensitive to the presence of buffer components and other contaminants, meaning MALDI is more robust.

You've just watched JoVE's video on MALDI mass spectrometry. This video described the theory behind the instrument, went over a general procedure, and covered some of the uses of the technique. Thanks for watching!

 Biochemistry

Soft Lithography

JoVE 5790

Many BioMEM devices, such as microfluidic channels, are fabricated using the soft lithography technique. Here, a microscale pattern is replicated by curing an elastomeric polymer over the 3D structure. These polymeric structures are then used to create a wide range of devices, ranging from microfluidic channels for biosensing applications to microscale bioreactors for the visualization of micro-colonies.

This video introduces photolithography and demonstrates the technique in the laboratory. Then, some applications of the technique and how the structures are used in the bioengineering field are examined.

 Bioengineering

Multi-Step Reactions

JoVE 11699

Chemical reactions often occur in a stepwise fashion involving two or more distinct reactions taking place in a sequence. A balanced equation indicates the reacting species and the product species, but it reveals no details about how the reaction occurs at the molecular level. The reaction mechanism (or reaction path) provides details regarding the precise, step-by-step process by which a reaction occurs. Each of the steps in a reaction mechanism is called an elementary reaction. These elementary reactions occur in sequence, as represented in the step equations, and they sum to yield the balanced chemical equation describing the overall reaction. In a multistep reaction mechanism, one of the elementary steps progresses slower than the others — sometimes significantly slower. This slowest step is called the rate-limiting step (or rate-determining step). A reaction cannot proceed faster than its slowest step, and hence, the rate-determining step limits the overall reaction rate.

Unlike balanced equations representing an overall reaction, the equations for elementary reactions are explicit representations of the chemical change. An elementary reaction equation depicts the actual reactant(s) undergoing bond-breaking/making and the product(s) formed. Rate laws may be derived directly from the balanced chemical equations for elementary reactions. However, this is not the case for most chemical reactions, where balanced equations often represent the overall change in the chemical system resulting from multistep reaction mechanisms. Therefore, the rate law must be determined from experimental data, and the reaction mechanism must be subsequently deduced from the rate law.

For instance, consider the reaction of NO2 and CO:

Figure1

The experimental rate law for this reaction at temperatures above 225 °C is:

Figure2

According to the rate law, the reaction is first-order with respect to NO2 and first-order with respect to CO. However, at temperatures below 225 °C, the reaction is described by a different rate law that is second-order with respect to NO2:

Figure3

This rate law is not consistent with the single-step mechanism, but it is consistent with the following two-step mechanism:

Figure4

Figure5

The rate-determining (slower) step gives a rate law showing second-order dependence on the NO2 concentration, and the sum of the two elementary equations gives the net overall reaction.

In general, when the rate-determining (slower) step is the first step in the reaction mechanism, the rate law for the overall reaction is the same as the rate law for this step. However, when the rate-determining step is preceded by an elementary step involving a rapidly reversible reaction, the rate law for the overall reaction may be more difficult to derive, often due to the presence of reaction intermediates.

In such instances, the concept that a reversible reaction is at equilibrium when the rates of the forward and reverse processes are equal can be utilized.

 Core: Organic Chemistry

Ethers from Alkenes: Alcohol Addition and Alkoxymercuration-Demercuration

JoVE 11740

Overview

Ethers can also be prepared from alkenes through acid-catalyzed addition of alcohols and alkoxymercuration–demercuration.

Preparation of Ethers by Acid-Catalyzed Addition of Alcohol to Alkenes

The acid-catalyzed addition of alcohol to an alkene involves treating the alkene with an excess of alcohol in the presence of an acid catalyst to form an ether under suitable conditions. The hydrogen will add to the less substituted carbon so that the nucleophile can attack the more substituted carbon across an alkene forming an ether.

Figure1

Preparation of Ethers by Alkoxymercuration–Demercuration

Alkoxymercuration–demercuration is a reaction in which an alkene and alcohol react in the presence of a mercuric acetate reagent followed by demercuration or reduction with sodium borohydride to yield an ether.

Figure2

The alkoxymercuration–demercuration mechanism follows Markovnikov's regioselectivity with the alkoxy group attached to the most substituted carbon and the H attached to the least substituted carbon. A variety of alcohols and alkenes can be used in the reaction. Ditertiary ethers cannot be prepared by this method due to steric hindrance.

 Core: Organic Chemistry

Regioselectivity and Stereochemistry of Acid-Catalyzed Hydration

JoVE 11777

The rate of acid-catalyzed hydration of alkenes depends on the alkene's structure, as the presence of alkyl substituents at the double bond can significantly influence the rate.

Figure1

The reaction proceeds with the slow protonation of an alkene by a hydronium ion to form a carbocation, which is the rate-determining step.

The reaction involving a tertiary carbocation intermediate is faster than a reaction proceeding through a secondary or primary carbocation. This can be justified by comparing their relative stabilities and the delocalization of the positive charge. Tertiary carbocations are the most stable and thus formed faster.

Regiochemical Outcome

The formation of a stable carbocation intermediate determines the regiochemical outcome as it directs the nucleophilic addition of water to the more substituted carbon following Markovnikov's orientation.

Stereochemical Outcome of Achiral Alkenes

For an achiral alkene such as 1-butene, protonation results in a secondary carbocation.

Figure2

The trivalent carbon is sp2-hybridized with a plane of symmetry. It can react with water either from the top or the bottom face with equal probability.

The reaction from the top face leads to (S)-2-butanol, while the reaction from the bottom face leads to (R)-2-butanol. Thus, the formation of a new chiral center leads to a racemic mixture of enantiomeric products.

Figure3

Stereochemical Outcome of a Chiral Alkene

The protonation of a chiral alkene forms a chiral carbocation with no plane of symmetry. The carbocation does not react equally from the top and bottom faces because one of the faces is more accessible than the other due to different steric setups, leading to a mixture of R and S products. Thus, two diastereomeric products are produced in unequal amounts, and the mixture is optically active.

Figure4

Rearrangement of a Carbocation

In some cases, the carbocation formed in the first step can rearrange to a more stable carbocation. For example, the protonation of 3-methyl-1-butene forms a 2° carbocation intermediate, which rearranges to a more stable 3° carbocation via a 1,2-hydride shift.

Figure5

 Core: Organic Chemistry

Structure and Physical Properties of Alkynes

JoVE 11831

Introduction:

In nature, compounds containing both carbon and hydrogen are known as "hydrocarbons". Aliphatic hydrocarbons are compounds whose molecules contain saturated single bonds (i.e., alkanes) or unsaturated double or triple bonds. Alkenes contain carbon–carbon double bonds and have a structural formula CnH2n. Unsaturated hydrocarbons containing carbon–carbon triple bonds are called "alkynes" and are structurally represented by the formula CnH2n-2.

The simplest alkyne is ethyne, or acetylene, a colorless gas, which undergoes combustion at high temperatures and is used as a fuel for welding. Alkynes are found in a variety of substances of both natural and synthetic origin. For example, natural alkynes can be found in the poison obtained from South American tree frogs. On the other hand, synthetic alkyne-containing compounds play an important role in drugs like oral contraceptives such as ethynylestradiol. Other alkyne-based drugs, such as selegiline, are used in combination with synthetic dopamine or L-dopa to treat Parkinson's disease. 

Hybridization of Alkynes:

The triple bond in alkynes occurs due to the overlap of the carbon–carbon (C–C) bond's sp orbitals forming a sigma (σ) bond, and the lateral overlapping of the 2py and 2pz orbitals forming the two pi (π) bonds, respectively. The overlap of the carbon sp orbital with the 1s orbital of the hydrogen atom results in the carbon–hydrogen (C–H) sp–1s sigma bond of the alkyne. Alkynes show an sp hybridization with a 50% s character. As electrons in the orbital show lower energies than the orbitals, an increase in the atom's s character increases its electronegativity.

Molecular Geometry: Bond length, Bond angles, and Bond Strengths

Owing to the sp hybridization involving the lateral overlap of the orbitals, alkynes have a linear geometry with the bond angle being 180°. The C–C triple bond's length in acetylene is measured to be 121 pm or 1.21 Å, which is lesser than that of alkenes (134 pm or 1.34 Å) and alkanes (153 pm or 1.53 Å) due to the increased number of bonds holding the two carbons together.

The C–H bond in alkynes (1.06 Å in acetylene) is also shorter compared to that in alkenes and alkanes due to the increased s-character of the sp hybridized carbon orbital forming the sp-1s σ bond with the hydrogen. In addition to having lower energies, electrons in the s-orbital are closer to the atomic nucleus and are bound more tightly than those in the p orbitals. As a result, a greater s-character increases the strength of the bond. Hence, C-C triple bonds and C-H bonds in alkynes are shorter and stronger than those of alkenes and alkanes, mainly due to their increased s character. The bond dissociation energy of alkynes is 966 kJ or 231 kcal/mol, which is higher than that of alkanes and alkenes due to the presence of these shorter and stronger bonds.

Physical Properties of Alkynes:

Alkynes show similar physical properties as their parent alkane or alkene. They are nonpolar compounds having a lower density than water and are insoluble in water and polar solvents. However, they show good solubility in nonpolar organic solvents. The lower-molecular weight alkynes such as ethyne and propyne exist as gases at room temperature, while higher molecular weight alkynes such as 1-octyne and 1-decyne are liquids.

 Core: Organic Chemistry

Microtubules in Cell Motility

JoVE 11914

Microtubules are thick hollow cylindrical proteins that help form the cytoskeleton. Microtubules have varied roles in the cell. These filaments help form cellular appendages like cilia and flagella, which are responsible for locomotion. The cilia arise from basal bodies, separated from the main body by a membrane-like structure forming the transition zone. This zone is the gate for the entry of lipids and proteins, creating a unique composition of lipids and proteins in the ciliary membrane and body. The central strand of the flagellum is called the axoneme and has a 9+2 arrangement of microtubules.

Microtubules help cells move using mechanisms like modulating actin polymerization by regulating Rho GTPase signaling pathways. During actin polymerization, with the help of +TIPs complex, microtubules sequester signaling molecules and actin assembly factors. These molecules are only released upon the disassembly of microtubules, thus regulating lamellipodia and filopodia formation.

Microtubules can also regulate directional migrations when they act as tracks for motor proteins to transport intracellular cargo and signaling molecules to the leading edge of the migrating cells. The cortical microtubules associated with the focal adhesion junctions help recycle focal adhesion proteins from the cell membrane during cell motility. They also facilitate the cross-talk between different cytoskeletal components. These microtubules undergo repeated cycles of rescue and catastrophe near the cell boundaries to regulate cell motility.

 Core: Cell Biology

Acid-Catalyzed Ring-Opening of Epoxides

JoVE 12109

Epoxides that are three-membered ring systems are more reactive than other cyclic and acyclic ethers. The high reactivity of epoxides originates from the strain present in the ring. This ring strain acts as a driving force for epoxides to undergo ring-opening reactions either with halogen acids or weak nucleophiles in the presence of mild acid. The acid catalyst converts the epoxide oxygen, a poor leaving group, into an oxonium ion, a better leaving group, making the reaction feasible. The reaction follows the SN2 mechanism, and the protonated oxygen, unlike other leaving groups, does not detach from the molecule. The regiochemistry of the products formed is governed by either steric or electronic effects. In the case of an asymmetrical epoxide bearing a primary and a secondary carbon, the steric effect dominates and favors the nucleophilic attack at a less-hindered carbon. However, in epoxides where one of the carbons is tertiary, the electronic effect comes into play and favors the attack at a more-hindered carbon. The stereochemistry of the products is similar to an SN2 reaction, where the nucleophile attacks anti to the leaving group. Notably, the anti-attack at a chiral carbon causes an inversion of configuration.

 Core: Organic Chemistry

Esters to β-Ketoesters: Claisen Condensation Overview

JoVE 12390

Regular Claisen condensation is a base-promoted reaction involving identical esters with two α hydrogens, condensing to produce β-ketoesters. It is a nucleophilic acyl substitution reaction wherein one of the ester molecules, upon deprotonation by the base, forms a nucleophilic enolate ion, while the other molecule serves as an electrophile.

Figure1

The condensation reaction requires specific nonaqueous bases, such as alkoxide, which must be similar to the alkoxy group of the ester molecule. The use of different alkoxide ions often leads to the transesterification product. Also, the use of hydroxide bases is avoided as they result in irreversible hydrolysis of the esters to carboxylate ions.

Figure2

 Core: Organic Chemistry

Regulation of Angiogenesis and Blood Supply

JoVE 12510

Rapidly dividing tumors, embryos, and wounded tissues require more oxygen than usual, lowering the oxygen concentration in the blood. At low oxygen or hypoxic conditions, an oxygen-sensitive transcription factor called the hypoxia-inducible factor 1 or HIF1 is activated. HIF1 is a dimeric protein of alpha (ɑ) and beta (β) subunits.  Under optimal oxygen conditions, HIF1β is present in the nucleus while HIF1ɑ remains in the cytosol. HIF1ɑ is hydroxylated by prolyl hydroxylase and factor inhibiting HIF-1. Hydroxylated HIF1 alpha binds von Hippel Lindau (VHL) E3 ubiquitin ligase and undergoes degradation. 

Hypoxia promotes HIF1ɑ accumulation. HIF1ɑ translocate to the nucleus and binds HIF1-beta, forming the HIF1 dimers. HIF dimer associates with CBP/P300 transcriptional regulator and binds HIF1 response elements of target genes, initiating their transcription. Some pro-angiogenic genes regulated by HIF1 include erythropoietin, angiopoietin, and vascular endothelial growth factor (VEGF). VEGF signaling is critical for regulating angiogenesis.

VEGF is a dimeric protein. They bind the transmembrane RTKs called  VEGF receptors (VEGFR). They have five isoforms, of which VEGF-A is most important in regulating angiogenesis and binds VEGFR1 with high affinity. VEGFA binding stimulates endothelial cells to differentiate into tip cells. Tip cells express high levels of delta-like notch ligand four or DLL4. DLL4 of tip cells interacts with the notch receptors of the neighboring cells promoting their differentiation to stalk cells.

Another essential ligand/ RTK signaling includes the angiopoietin/ tie-2 receptor that works closely with the VEGF signaling pathway in the latter stages of angiogenesis.  Angiogepoeitin 1/ tie-2 signaling promotes endothelial cell survival, initiates vascular branching, and helps stabilize newly formed vessels.

A third ephrin-B/ephrin-B RTK signaling pathway is also essential for angiogenesis and helps specify endothelial cells into arterial and venous cell types.

Apart from these ligand/receptor interactions, some molecules help mediate cell-cell interactions and cell-matrix interactions and regulate angiogenesis. For example, matrix metalloproteases allow basement membrane degradation and endothelial tip migration to the target tissue. Alternatively, protease inhibitors prevent matrix degradation and stabilize them once the vessel is formed. VE-Cadherins, N cadherins, and occludin are essential junctional proteins that stabilize the newly formed endothelial lining.

 Core: Cell Biology

Alkylation of β-Ketoester Enolates: Acetoacetic Ester Synthesis

JoVE 13074

Acetoacetic ester synthesis is a method to obtain ketones from alkyl halides and β-keto esters. The reaction occurs in the presence of an alkoxide base that abstracts the acidic proton of the β-keto esters. The step results in an enolate ion which is doubly stabilized. The enolate then reacts with an alkyl halide via the SN2 process to produce an alkylated ester intermediate with a new C–C bond. The hydrolysis of the intermediate, followed by acidification, results in an alkylated β-keto acid. Under high-temperature conditions, the β-keto acid undergoes decarboxylation to form a ketone. However, if the alkylation is repeated before hydrolysis and decarboxylation steps, a disubstituted ketone is obtained.

 Core: Organic Chemistry

Subcellular Fractionation

JoVE 13375

The homogenate obtained after cell lysis contains various membrane-bound organelles that can be further separated into pure fractions by subcellular fractionation. These isolates are used to study specific cellular components, analyze localized protein activity, and are even employed in diagnostics. Fractionation is typically achieved using centrifugation methods, the most common being density-gradient and differential centrifugation.

Differential Centrifugation

Differential centrifugation is a relatively simple  method that separates the cellular components based on size and density. Sequential centrifugation with increasing speeds (ranging from 10,000 X g to 150,000 x g) sediments the differently sized components. However, since multiple organelles can be of similar size and density, this method usually produces crude fractions.

Density Gradient Centrifugation:

Highly purified fractions of cellular components can be obtained by separating the homogenate in a density gradient solution. A density gradient is prepared in a centrifuge tube by layering solutions with increasing densities, such as increasingly concentrated sucrose solutions, with the densest layer at the bottom of the tube. Such gradients are used in rate-zonal centrifugation to separate cellular organelles based on their size and shape. Upon centrifugation, the organelles sediment at different rates, based on their sedimentation coefficients, as they move through the different density layers.

Alternatively, a continuous density gradient can also be prepared by mixing solutions of different densities in gradual proportions along the length of the tube. During centrifugation, each component immobilizes at the position that matches their density — their equilibrium position. Hence this method is also known as equilibrium or buoyant sedimentation. This separation of cellular components and molecules is thus based on their density, not their size.

 Core: Cell Biology

Fixation and Sectioning

JoVE 13391

Two basic types of preparation are used to visualize specimens with a light microscope: wet mounts and fixed specimens.

The simplest type of preparation is the wet mount, in which the specimen is placed in a drop of liquid on the slide. A liquid specimen can be directly deposited on the slide using a dropper. Solid specimens, such as skin scraping, can be placed on the slide before adding a drop of liquid to prepare the wet mount. Sometimes the liquid is simply water, but stains are often added to enhance contrast. Once the liquid has been added to the slide, a coverslip is placed on top, and the specimen is ready for examination under the microscope.

The second method of preparing specimens for light microscopy is fixation. The “fixing” of a sample refers to the process of attaching cells to a slide. Fixation is often achieved by heating (heat fixing) or chemically treating the specimen. In addition to attaching the specimen to the slide, fixation also kills microorganisms in the specimen, stopping their movement and metabolism while preserving the integrity of their cellular components for observation.

To heat-fix a sample, a thin layer of the specimen is spread on the slide (called a smear), and the slide is then briefly heated over a heat source. Chemical fixatives are often preferable to heat for tissue specimens. Chemical agents such as acetic acid, ethanol, methanol, formaldehyde (formalin), and glutaraldehyde can denature proteins, stop biochemical reactions, and stabilize cell structures in tissue samples.

In addition to fixation, staining is almost always applied to color certain features of a specimen before examining it under a light microscope. Stains, or dyes, contain salts made up of a positive ion and a negative ion. Depending on the type of dye, the positive or the negative ion may be the chromophore (the colored ion); and the other uncolored ion is called the counterion. If the chromophore is the positively charged ion, the stain is classified as a basic dye; if the negative ion is the chromophore, the stain is considered an acidic dye.

Dyes are selected for staining based on the chemical properties of the dye and the specimen being observed, which determine how the dye will interact with the specimen. In most cases, it is preferable to use a positive stain, a dye that will be absorbed by the cells or organisms being observed, adding color to objects of interest to make them stand out against the background. However, there are scenarios where it is advantageous to use a negative stain absorbed by the background but not by the cells or organisms in the specimen. Negative staining produces an outline or silhouette of the organisms against a colorful background.

Because cells typically have negatively charged cell walls, the positive chromophores in basic dyes tend to stick to the cell walls, making them positive stains. Thus, basic dyes such as basic fuchsin, crystal violet, malachite green, methylene blue, and safranin typically serve as positive stains. On the other hand, the negatively charged chromophores in acidic dyes are repelled by negatively charged cell walls, making them negative stains. Commonly used acidic dyes include acid fuchsin, eosin, and rose bengal.

Some staining techniques involve the application of only one dye to the sample; others require more than one dye. In simple staining, a single dye is used to emphasize particular structures in the specimen. A simple stain will generally make all of the organisms in a sample appear the same color, even if the sample contains more than one type of organism. In contrast, differential staining distinguishes organisms based on their interactions with multiple stains. In other words, two organisms in a differentially stained sample may appear different colors. Differential staining techniques commonly used in clinical settings include Gram staining, acid-fast staining, endospore staining, flagella staining, and capsule staining.

This text is adapted from Openstax, Microbiology 2e, Section 2.4: Staining microscopic specimens

 Core: Cell Biology

Interphase

JoVE 13407

The cell cycle occurs over approximately 24 hours (in a typical human cell) and in two distinct stages: interphase, which includes three phases of the cell cycle (G1, S, and G2), and mitosis (M). During interphase, which takes up about 95 percent of the duration of the eukaryotic cell cycle, cells grow and replicate their DNA in preparation for mitosis.

Phases of Interphase

Following each period of mitosis and cytokinesis, eukaryotic cells enter interphase, during which they grow and replicate their DNA in preparation for the next mitotic division.

During the G1 (gap 1) phase, cells grow continuously and prepare for DNA replication. During this phase, cells are metabolically active and copy essential organelles and biochemical molecules, such as proteins.

In the subsequent S (synthesis) phase of interphase, cells duplicate their nuclear DNA, which remains packaged in semi-condensed chromatin. During the S phase, cells also duplicate the centrosome, a microtubule-organizing structure that forms the mitotic spindle apparatus. The mitotic spindle separates chromosomes during mitosis.

In the G2 (gap 2) phase, which follows DNA synthesis, cells continue to grow and synthesize proteins and organelles to prepare for mitosis.

In human cells, the G1 phase spans approximately 11 hours, the S phase takes about 8 hours, and the G2 phase lasts about 4 hours. During G1, cells are diploid (2n, a pair of each chromosome). Following replication in S phase, cells increase their DNA content to 4n. Cells remain 4n until cytokinesis, at which point their DNA content is reduced to 2n.

 Core: Cell Biology

Nondisjunction

JoVE 13439

Nondisjunction is the failure of homologous chromosomes or sister chromatids to separate correctly and move to the opposite poles of the cells. This produces daughter cells with abnormal chromosome numbers.  Nondisjunction is common during anaphase I or anaphase II of meiosis.  Mutations in synaptonemal complex proteins that attach homologous chromosomes increase the chances of nondisjunction in anaphase I of meiosis I. In contrast, mutations in topoisomerases and condensins that hold sister chromatids together promote nondisjunction during anaphase II of meiosis II.

Nondisjunction of chromosomes in germ cells results in gametes possessing additional or fewer chromosomes than normal. Nondisjunction is more frequent during oogenesis than during spermatogenesis. When a gamete with abnormal chromosomes fertilizes a gamete with a normal chromosome number, the resulting zygote has an abnormal number of chromosomes or aneuploidy. Such aneuploid zygotes can have fewer chromosomes than normal, leading to monosomy (45; 2n-1), or more chromosomes than normal, leading to trisomy (47; 2n+1). Some females lack one of the X chromosomes, a typical case of monosomy (45, X), and develop Turner Syndrome. In other instances, individuals who develop Down Syndrome have trisomy with three copies of chromosome 21.

 Core: Cell Biology

Tissue Renewal without Stem Cells

JoVE 13471

After cellular or tissue damage, the resident stem cells present in the human body can locally repair and regenerate the damaged tissue or organ. However, even though some tissues do not have stem cells, they can repair and regenerate with the help of pre-existing cells. For example, beta cells of the pancreas and hepatocytes of the liver can divide to renew and regenerate the tissue. Here, both cell division and cell death are well regulated by homeostasis.

However, failure of such a system may result in life threatening diseases. Dysfunction in insulin-secreting cells, as well as the target cells’ responsiveness to insulin, can lead to a condition called diabetes mellitus. Interestingly, a recent study also found that as a backup, both the liver and pancreas have a few stem cells which are activated under extreme conditions to produce differentiated cell types. In this case, both the liver and pancreas revert back to their normal mechanisms of repair and renewal.

 Core: Cell Biology

Types Of Column Chromatography

JoVE 13560

The stability and compatibility of column material with samples are crucial for efficient purification in chromatographic techniques. Various operating parameters such as pH, temperature, or solvent affect the packing of the column material, thereby determining the purification efficiency. The choice of column material also plays an essential role in deciding the operating parameters and can be modified based on the proteins that need to be purified.

Gel Filtration Chromatography

When the protein's chemical nature is unknown, gel filtration chromatography is used for purification based on size. Matrices such as agarose, polyacrylamide, dextran derivatives, or silica beads are used for low-pressure systems, whereas polymeric resins are used for higher flow rates. The varying degree of cross-linking of these polymeric materials determines the pore size of the matrix, which decides the size range of proteins that can be separated.

Ion Exchange Chromatography

Ion-exchange chromatography can purify proteins based on their net charge at a particular pH determined by their isoelectric point or pI. The net charge on matrix resin is used to purify oppositely charged proteins by replacing the cations or anions from the matrix resin with the charged proteins. The resins are categorized based on the type of ion that is exchanged — positively charged proteins replace cations from cation-exchange resins, and vice versa. Strong ion exchangers with stability over a wide range of buffer pH include quaternary ammonium, sulfonate, and sulfopropyl resins. In contrast, weak ion exchangers that are effective only in a narrow pH range include diethyl aminoethyl (DEAE) or carboxymethyl (CM) resins.

Affinity Chromatography

Affinity chromatography involves immobilizing protein-specific antibodies or enzyme-specific substrates to the column material, usually agarose beads. For example, a column with the immobilized antibody of choice is prepared to separate specific antigens. When the mixed antigen sample is passed through the  column, the target antigens bind to the immobilized antibodies and separate from the mixture. To separate recombinant proteins tagged with enzymes such as beta-galactosidase, the corresponding substrate Isopropyl-β-D-thiogalactopyranoside (IPTG) is immobilized on the column to purify the enzyme.

High-pressure liquid chromatography (HPLC) and gas chromatography (GC) are the most commonly used column chromatography techniques. In HPLC, the solvent flows at high pressure through a tightly packed column that improves efficiency. Gas chromatography analyzes volatile compounds by converting them to their gaseous forms upon heating. The volatilized sample is pushed through the column by inert gasses such as nitrogen and helium.

 Core: Cell Biology

Statistical Analysis: Overview

JoVE 14505

When we take repeated measurements on the same or replicated samples, we will observe inconsistencies in the magnitude. These inconsistencies are called errors. To categorize and characterize these results and their errors, the researcher can use statistical analysis to determine the quality of the measurements and/or suitability of the methods.

One of the most commonly used statistical quantifiers is the mean, which is the ratio between the sum of the numerical values of all results and the total number of results. Another commonly used quantifier is the median, which is the middle value amongst all the results arranged in the increasing or decreasing order of their numerical values. In the presence of extremely large or small values in the data set, the median can be a more suitable measure for the data set overall. Both the mean and median can be useful measures of the central value of a set of measurements. In addition to the central value, the range is another way to characterize the distribution of a set of values. It is the numerical difference between the highest and lowest values in a set of results.

Precision and accuracy are two important measures of observed errors in the data set. Precision is the measure of closeness between the replicate measurements. In a precise data set, the values of different measurements cluster close together. The values are farther apart or more scattered in an imprecise data set. Often, the range of an imprecise data set will be high. On the other hand, accuracy is a measure of closeness between the measurements and the true or expected value. This means that the magnitude of errors in a more accurate data set is smaller than that of a less accurate data set.

 Core: Analytical Chemistry

Difference from Background: Limit of Detection

JoVE 14521

The limit of detection (LOD) is the smallest amount of analyte that can be distinguished from the background noise. The LOD value corresponds to the concentration at which the analyte signal is three times larger than the standard deviation of the blank signal. Below this value, the analyte signal cannot be differentiated from the background noise. It is calculated by dividing the calibration slope by 3 times the standard deviation of the blank signals.

The LOD indicates the presence or absence of an analyte but is usually too low to be reliably quantified. For quantification, we need another value called the limit of quantification, which is defined as the lowest quantity of analyte that the instrument can quantify. Its value corresponds to the concentration at which the signal is ten times larger than the standard deviation of the blank signal.

 Core: Analytical Chemistry

Complexation Equilibria: The Chelate Effect

JoVE 14537

In complexation reactions, metal atoms or cations interact with ligands to form donor-acceptor adducts called metal complexes. Ligands that bind through one donor site are monodentate, ligands with two donor sites are bidentate, and those with more than two donor sites are polydentate ligands. For example, ethylene diamine is a bidentate ligand that binds through two nitrogen donor atoms, forming a five-membered ring. EDTA is a polydentate ligand that binds through four oxygen and two nitrogen atoms.

Bidentate and polydentate ligands are also called chelating agents, and the corresponding complexes are called chelates. Chelate is a Greek word that means "claw-like." Metal complexes formed by the chelating agents are more stable than those formed by their monodentate counterparts, as the reaction for their formation is entropically favored. This property is known as the chelate effect or the entropy effect.

 Core: Analytical Chemistry

Complexometric EDTA Titration Curves

JoVE 14576

EDTA titration curves determine the free metal ion concentration. The titration curve represents the change in concentration of free metal ions (p function) as a function of the volume of EDTA added. This curve consists of three regions: before, at, and after equivalence points. Excess free metal ions are present before the equivalence point. Equal concentrations of metal ions and EDTA are present at the equivalence point. After the equivalence point, excess EDTA exists. This means slight dissociation can be observed at and after the equivalence point.

The complex's conditional formation constant (Kf′) calculates the free metal ion concentration at and after the equivalence point, and the shape of the titration curve is affected by Kf′ of the complex. For example, the Ca–EDTA complex has a larger Kf′ than the Sr–EDTA complex. As a result, the Ca–EDTA titration curve has a larger break at the equivalence point.

The Kf′ of the complex depends on the pH of the solution. For instance, Ca–EDTA exhibits various shapes at different pH. At higher pH, Ca–EDTA has a larger Kf′, and complex formation is more favorable. The curve has a large break at the equivalence point. At lower pH, the Kf′ of the complex is small, indicating less favorable complex formation. As a result, the curve has a small break at the equivalence point.

 Core: Analytical Chemistry

Mass Spectrometry: Complex Analysis

JoVE 14592

Mass spectrometry is an important technique for the identification of pure compounds. However, it has some limitations for the analysis of complex mixtures, often due to excessive fragmentation making the spectrum too complicated to decipher. Mass spectrometry can be combined with suitable separation methods in sequence, forming hyphenated methods, which are useful in the analysis of complex mixtures.

GC–MS is a powerful hyphenated method commonly used in forensics and environmental laboratories for precise analysis of mixtures. Gas chromatography uses narrow capillary columns to separate components of a thermally stable volatile mixture and passes them to the mass spectrometer for analysis.

LC–MS is another hyphenated method that couples liquid chromatography with a mass spectrometer to analyze nonvolatile mixtures. To make liquid chromatography compatible with a mass spectrometer, suitable pressure-maintaining ionization interfaces or atmospheric-pressure ionization techniques like electrospray ionization, which is applicable for polar and ionic compounds, or atmospheric-pressure chemical ionization, which applies to less polar molecules, are used.

Another hybrid method, called tandem mass spectroscopy, uses multiple mass analyzers in sequence. Compared to other hyphenated methods, tandem spectroscopy is faster, more sensitive, and more selective due to smaller chemical noise. Tandem spectroscopy can be further combined with separation techniques to form GC–MS/MS or LC–MS/MS for more complex mixture analysis.

Capillary electrophoresis–mass spectrometry is a very sensitive technique commonly used to analyze large biomolecules like DNA, proteins, and polypeptides. This method feeds quadrupole mass analyzers with capillary effluents after passing through an electrospray ionization interface.

 Core: Analytical Chemistry

Ion Exchange

JoVE 14625

Ion exchange chromatography separates charged molecules from a solution by reversibly exchanging them with mobile, or 'active', ions associated with the oppositely charged stationary phase. This method can be used to separate ions, soften and deionize water, and purify solutions. The polymers comprising the ion-exchange column are high-molecular-weight and chemically stable polymers, crosslinked to be porous and essentially insoluble. They are also functionalized with either acidic or basic groups. When functionalized with acidic groups, the columns or resin are known as cation exchangers and contain negatively charged polymers with positive counter ions. These allow cations to be exchanged from the mobile phase. When functionalized with basic groups, the columns or resin are known as anion exchangers, which contain positively charged polymers with negative counter ions and exchange anions from the solution.

During ion exchange, the counterions associated with the solid-phase polymers enter the solution while ions from the solution interact with the polymer in their place. Note that the ability to exchange ions increases as the extent of cross-linking increases and the number of ion-exchange groups increases in the polymer. This exchange is reversible and continues until equilibrium is established.

Acid cation exchangers can be functionalized with either strong acids, such as sulfonic acid (−SO3H), seen in polystyrene sulfonic acid, or weak acids, such as carboxylic acid (−COOH), seen in poly(methyl methylacrylic) acid. Similarly, anion exchangers are classified as either strong or weak. Strong anion exchangers often contain quaternary ammonium groups (−CH2N(CH3)3+), such as those in polystyrene quaternary ammonium chloride, and weak anion exchangers are often functionalized with substituted amines (−NH3+), such as those in polystyrene tertiary amine hydroxide.

 Core: Analytical Chemistry

Specialized Characteristics of Cardiac Muscles

JoVE 14854

The primary role of cardiac muscles is to propel blood throughout the cardiovascular system. The cardiac muscle cells, or cardiomyocytes, exhibit specialized characteristics that allow them to perform this function.

Cardiac muscle cells are smaller than skeletal muscles, averaging 10–20 mm in diameter and 50–100 mm in length. However, they have large energy demands for continuous contraction and relaxation. This energy is almost exclusively derived from aerobic metabolism of energy reserves in the form of glycogen and lipid inclusions. Additionally, the sarcoplasm of these cells contains large numbers of mitochondria and abundant reserves of myoglobin, which store the oxygen needed to break down energy reserves during peak activity.

In contrast to skeletal muscles, cardiomyocytes can contract independently without relying solely on nerve stimulation. The pacemaker cells, a distinct type of cardiomyocyte, can generate action potentials and initiate contractions in cardiac muscles approximately 75 times per minute. This depolarization wave swiftly spreads through the muscle cells through gap junctions, prompting all the cells to contract and relax together as one unit, known as the functional syncytium.

Cardiomyocytes also differ from skeletal muscles in that the depolarization wave triggers the gradual release of calcium ions into the sarcoplasm of the contractile cells from both the sarcoplasmic reticulum and the interstitial fluid. This results in the contraction cycle of cardiac muscles that lasts for about 200 ms, which is longer than the skeletal muscle contractions that last only 40 to 120 ms. The extended cycle also increases the absolute refractory period of cardiac muscles, which coincides with the repolarization phase of the action potential. This feature enables the muscle to fully relax before another action potential is generated, preventing sustained or tetanic contractions.

 Core: Anatomy and Physiology

Muscles of the Shoulder

JoVE 14875

The muscles surrounding the shoulder girdle, including the clavicle and scapula, primarily stabilize the scapula. This stable base allows other muscles to move the humerus effectively. Scapular movements often mirror those of the humerus and extend its range of motion. For instance, raising the arm above the head would not be feasible without simultaneous upward rotation of the scapula.

Anterior Thoracic Muscles

The anterior thoracic muscles include the serratus anterior, subclavius, and pectoralis minor. The serratus anterior, often known as the boxer’s muscle, starts from the first to eighth or ninth ribs. It inserts along the superior and inferior angle of the scapula as well as its medial border. This muscle is crucial for abducting and rotating the scapula upward. When the scapula is fixed, it helps elevate the ribs, assisting in punching movements.

The pectoralis major joins the pectoral girdle to the thorax. It originates from the anterior surface of the superior margins of the second or third to fourth or fifth ribs and inserts on the coracoid process of the scapula. It tilts the scapula anteriorly and pulls it inferiorly. It also elevates the thorax, aiding in respiration.

The subclavius muscle originates on the first rib and inserts into the lower surface of the lateral clavicle. Its primary function is to stabilize the clavicle during shoulder and arm movements. It also depresses the clavicle and moves it anteriorly.

Posterior Thoracic Muscles

The posterior thoracic muscles include the trapezius, levator scapulae, and the rhomboid minor and major muscles. The trapezius muscle, one of the broadest and most distinct muscles, has its roots in the bony structure of the cervical and thoracic regions. The muscle originates from the external protuberance of the occipital bone, ligamentum nuchae, and the spinous process of the seventh cervical and all thoracic vertebra. It inserts at the clavicle along its lateral third in addition to the acromion and the spine of the scapula. It is categorized into three parts — the upper, middle, and lower fibers, each performing unique actions. The upper fibers elevate the scapula and extend the neck, the middle fibers adduct the scapula, and the lower fibers depress the scapula.

The levator scapulae originate from the transverse process of the cervical vertebrae, C1-C4, and insert into the scapulae at the superior part of its medial border. It elevates the scapula and causes its downward rotation. It also flexes the neck laterally.

Originating from the nuchal ligament and spinous processes of C7 and T1, the rhomboid minor inserts into the vertebral border of the scapula, superior to the spine. On the other hand, the rhomboid major originates from the spinous processes of the T2 to the vertebral border of the scapula, inferior to the spine. These muscles help retract the scapula and rotate it medially.

 Core: Anatomy and Physiology

Association Areas of the Cortex

JoVE 14908

Association areas are regions of the cerebral cortex that do not have a specific sensory or motor function. Instead, they integrate and interpret information from various sources to enable higher cognitive processes such as memory, learning, and decision-making. Some key association areas include the following:

Prefrontal Association Area: This area is located in the frontal lobe and is involved in planning, decision-making, and moderating social behavior. It connects with primary motor areas, sensory areas, and the limbic system, which processes emotions.

Parietotemporal Association Area: Found at the junction of the parietal and temporal lobes, this area is crucial for understanding language and spatial awareness. It receives input from the primary auditory cortex, somatosensory cortex, and other sensory areas.

Limbic Association Area: This area is essential for memory formation and emotional processing. Located in the medial part of the temporal lobe, it connects with the hippocampus, amygdala, and other parts of the limbic system.

Broca's Area: Broca's area, located in the posterior part of the left inferior frontal gyrus, plays a crucial role in speech production. It connects with the primary motor cortex, responsible for controlling speech muscles, and Wernicke's area through a bundle of nerve fibers called the arcuate fasciculus.

Damage to Broca's area can result in Broca's aphasia, a condition characterized by non-fluent speech, difficulty forming complete sentences, and trouble finding the right words. However, comprehension usually remains relatively intact.

Wernicke's Area: Wernicke's area, situated in the superior temporal gyrus of the left temporal lobe, is responsible for language comprehension. It receives input from primary auditory and visual cortices, allowing it to process spoken and written language. Wernicke's area connects with Broca's through the arcuate fasciculus, enabling communication between speech production and comprehension.

Damage to Wernicke's area can lead to Wernicke's aphasia, a condition where individuals have fluent but nonsensical speech with poor comprehension. They may also need help understanding written language.

Association Areas of Special Senses: The association areas of special senses are responsible for processing and interpreting information from our primary sensory cortices. These include:

Visual Association Area: The visual association area, or the visual association cortex, is located in the brain's occipital lobe. This region plays a crucial role in processing and interpreting visual information, allowing us to make sense of what we see. It receives input from the primary visual cortex (V1), which detects basic features such as edges, colors, and motion. The visual association area then integrates this information to create a more comprehensive understanding of the visual scene, enabling us to recognize objects, faces, and other complex elements. Damage to the visual association area can result in visual agnosia, where individuals can see objects but cannot recognize or identify them.

Auditory Association Area: Found in the temporal lobe, this area interprets sounds and speech from the primary auditory cortex. Damage to the auditory association area can lead to auditory agnosia, where individuals can hear sounds but cannot understand or interpret them.

Somatosensory Association Area: This area is situated in the parietal lobe. It integrates and interprets tactile information from the primary somatosensory cortex. Damage to the somatosensory association area can result in difficulties recognizing objects by touch, astereognosis, or identifying body parts' position in space, and agnosia for body schema.

Orbitofrontal Cortex

The orbitofrontal cortex (OFC) is a region in the brain's frontal lobes, situated just above the orbits of the eyes. It involves various cognitive functions, including decision-making, emotional regulation, and reward processing. While not directly responsible for visual processing, the OFC is interconnected with other brain regions that handle visual information, such as the visual association cortex and the fusiform face area.

The OFC contributes to our ability to make decisions based on visual stimuli, such as choosing between different food items or evaluating the attractiveness of a potential mate. Additionally, it helps process the emotional content of visual information, allowing us to respond appropriately to facial expressions and other emotionally charged visual cues.

Facial Recognition Area

The facial recognition area, commonly known as the fusiform face area (FFA), is located in the fusiform gyrus of the temporal lobe. This region is dedicated explicitly to recognizing faces and distinguishing them from other objects. The FFA allows us to identify familiar faces quickly and accurately interpret facial expressions and emotions. The FFA works with other brain regions, such as the occipital face area (OFA) and the superior temporal sulcus (STS), to create a robust facial recognition system. While the FFA focuses on processing the invariant aspects of faces (e.g., identity), the OFA and STS deal with more changeable aspects, such as gaze direction and expression.

Frontal Eye Field Area

The frontal eye field (FEF) is a region within the frontal cortex that explicitly controls eye movements and visual attention. The FEF guides our eyes toward relevant visual stimuli and away from irrelevant or distracting information. It coordinates saccadic eye movements, which are rapid, voluntary shifts in gaze that allow us to explore our visual environment efficiently. Additionally, the FEF contributes to higher-order cognitive processes, such as working memory and attention. The FEF enables us to prioritize and process relevant information more effectively by directing our visual attention to specific locations or objects.

 Core: Anatomy and Physiology

Spinal Nerves: Anatomy

JoVE 14924

Spinal nerves are pivotal conduits in the nervous system, bridging the central nervous system (CNS) with the peripheral nervous system (PNS). These nerves enable a complex communication network between the brain, spinal cord, and the rest of the body, facilitating sensory input, motor output, and autonomic functions.

There are 31 bilateral pairs of spinal nerves, each emerging from the spinal cord through the intervertebral foramina—openings between adjacent vertebrae. These nerves are symmetrically organized and are categorized based on the region of the spine from which they emerge —  8 cervical (C1-C8), 12 thoracic (T1-T12), 5 lumbar (L1-L5), 5 sacral (S1-S5), and 1 coccygeal (Co1).

  • • Cervical nerves (C1-C8): These eight nerve pairs transmit signals to and from the head, neck, and upper extremities. The cervical nerves are:
    • C1: Suboccipital nerve
    • C2: Greater occipital nerve
    • C3: Lesser occipital nerve
    • C4: Supraclavicular nerve
    • C5: Dorsal scapular nerve
    • C6: Superior plexus root
    • C7: Middle plexus root
    • C8: Inferior plexus root
  • • Thoracic nerves (T1-T12): The 12 pairs of thoracic nerves correspond to the 12 thoracic vertebrae and control the sensation and motor function of the chest, upper back, and abdominal muscles. The T1 to T11 nerves are the intercostal nerves, while T12 is the subcostal nerve.
  • • Lumbar nerves (L1-L5): Five pairs of lumbar nerves control the lower back, hips, and parts of the legs. The five lumbar nerves are:
    • L1: Iliohypogastric nerve
    • L2: Ilioinguinal nerve
    • L3: Genitofemoral nerve
    • L4: Lateral femoral cutaneous nerve
    • L5: Lumbar plexus root
  • • Sacral nerves (S1-S5): The sacral nerves consist of five pairs, and they are responsible for the function of the lower extremities, pelvic organs, and some muscles in the hips and legs. The sacral nerves are:
    • S1: Superior gluteal nerve
    • S2: Inferior gluteal nerve
    • S3: Sciatic nerve
    • S4: Pudendal nerve
    • S5: Sacral plexus root
  • • Coccygeal nerve (Co1): The single coccygeal nerve pair, called anococcygeal nerve, provides sensory innervation to the skin over the coccyx and surrounding area, as well as to some of the muscles in the pelvic floor. Its motor functions are minimal and not as well defined as those of the other spinal nerves.

Structure of Spinal Nerve

Each spinal nerve is formed by the union of two distinct roots — the dorsal (posterior) root and the ventral (anterior) root.

  • • Dorsal Root: The dorsal root carries sensory information from the peripheral receptors to the spinal cord. It contains sensory neurons whose cell bodies are located in the dorsal root ganglion, a bulge in the root just before it merges with the ventral root.
  • • Ventral Root: In contrast, the ventral root conveys motor information from the spinal cord to the muscles and glands. This root consists of axons from motor neurons whose cell bodies are housed within the gray matter of the spinal cord.

These two roots combine to form a mixed spinal nerve carrying both sensory and motor fibers. Shortly after emerging from the spinal column, each spinal nerve splits into two primary branches — the dorsal ramus and the ventral ramus.

  • • Dorsal Ramus: This branch innervates the muscles and skin of the posterior trunk, providing sensory and motor functions to these areas.
  • • Ventral Ramus: The larger of the two branches, the ventral ramus, supplies the anterolateral parts of the trunk and the limbs. It carries fibers that form the nerve plexuses — networks of interconnecting nerves that give rise to peripheral nerves innervating the limbs.

Additionally, each spinal nerve has a small meningeal branch that reenters the vertebral column to innervate the vertebrae, ligaments, blood vessels, and meninges of the spinal cord. Spinal nerves also feature communicating rami that connect with the sympathetic trunk of the autonomic nervous system (ANS). These rami consist of:

  • • Gray Rami Communicantes: These carry postganglionic sympathetic fibers back to the spinal nerves for distribution to the target organs and tissues.
  • • White Rami Communicantes: These are present only in the thoracic and upper lumbar regions. They carry preganglionic sympathetic fibers from the spinal nerve to the sympathetic trunk.

 Core: Anatomy and Physiology

Sleep-Wake Cycles

JoVE 14943

Sleep is an essential physiological process vital to maintaining overall well-being. The reticular activating system (RAS), a network of neurons in the brainstem, regulates wakefulness and sleep. While it may seem passive, sleep consists of distinct cycles, each with its unique characteristics and functions. Two key sleep phases are non-rapid eye movement (NREM) and  rapid eye movement (REM).

NREM Sleep

NREM sleep comprises four progressive stages that seamlessly merge:

  1. Stage 1  is the transitional phase bridging wakefulness and sleep. This stage usually lasts between 1 and 7 minutes, during which the person feels relaxed and has closed eyes. Individuals awakened during this stage often express a sense of not having been asleep.
  2. Stage 2, characterized as light sleep, marks the onset of the sleeping state. It is relatively easy to rouse someone during this stage. Fragmented dreams may occur, and the eyes may exhibit slow horizontal movements.
  3. Stage 3 denotes a period of moderately deep sleep. The individual experiences decreased body temperature and blood pressure, making it slightly more challenging to awaken them. Typically, Stage 3 emerges about 20 minutes after falling asleep.
  4. Stage 4 signifies the deepest level of sleep. Although brain metabolism considerably decreases and there is a slight drop in body temperature, most reflexes remain intact, and muscle tone only undergoes minimal reduction. Waking an individual during this stage becomes exceedingly difficult.

REM Sleep

REM sleep is an intriguing phase characterized by rapid eye movements and intense brain activity including vivid dreams. While the body appears still and relaxed, internally, the brain is highly active. The muscles are temporarily paralyzed to prevent acting out dreams.

Sleep Disorders

Sleep-related health issues are a serious concern and can manifest in various forms, including insomnia, sleep apnea, and narcolepsy. Insomnia is characterized by persistent difficulty in falling or maintaining sleep, often as a consequence of stress, overconsumption of caffeine, disrupted sleep-wake cycles, or depression. Sleep apnea, in contrast, is a condition where the sufferer's breathing involuntarily stops for periods exceeding 10 seconds during sleep, usually due to the throat's muscular tone diminishing, causing the airway to constrict. Narcolepsy is unique in that it is not associated with an inability to sleep but rather a failure to stay awake. Sufferers experience sudden bouts of sleep during their waking hours. Research has linked this to a lack of the neuropeptide orexin, also known as hypocretin, which aids in maintaining wakefulness. A malfunction in the hypothalamus, where orexin is produced, can be a critical factor in the development of narcolepsy.

 Core: Anatomy and Physiology

Anatomy of the Eyeball

JoVE 14962

The eye is a spherical, hollow structure composed of three tissue layers. The outer layer — the fibrous tunic, comprises the sclera — a white structure — and the cornea, which is transparent. The sclera encompasses some of the ocular surface, most of which is not visible. However, the 'white of the eye' is distinctively visible in humans compared to other species. The cornea, a clear covering at the front of the eye, enables light penetration. The eye's middle layer, the vascular tunic, primarily consists of the choroid, the ciliary body, and the iris. The choroid is a highly vascularized connective tissue supplying blood to the eyeball, situated behind the ciliary body. The ciliary body, a muscular entity, is linked to the lens by zonule fibers or suspensory ligaments. These aid in lens curvature, facilitating the focus of light onto the rear of the eye. The iris, the eye's colored portion, overlays the ciliary body and is visible at the front of the eye. Iris, a circular muscle, dilates or constricts the pupil, the central eye aperture that permits light entry. The iris contracts the pupil in bright light, which widens the pupil in dim light. The innermost layer, the neural tunic or retina, houses the nervous tissue in light perception.

The eye can be segmented into two distinct sections: the front cavity and the back cavity. The front cavity between the cornea and the lens — encapsulating the iris and ciliary body — is filled with a light liquid known as aqueous humor. On the other hand, the back cavity expands from the area behind the lens to the back of the inner eyeball, where the retina is positioned. This cavity is filled with a thicker fluid referred to as vitreous humor.

The retina is a complex structure composed of numerous layers with distinct cells dedicated to the preliminary processing of visual signals. Photoreceptors, namely rods and cones, respond to light energy by altering their membrane potential. This change influences the quantity of neurotransmitters that the photoreceptors dispatch onto the bipolar cells in the outer synaptic stratum. In the retina, it is the bipolar cell that links a photoreceptor to a retinal ganglion cell (RGC) situated in the inner synaptic layer. Amacrine cells aid in processing within the retina before the RGC generates an action potential. Positioned at the retina's deepest layer, the RGCs' axons aggregate at the optic disc and exit the eye, forming the optic nerve. Since these axons traverse the retina, there is an absence of photoreceptors at the eye's rear, where the onset of the optic nerve lies. This results in a "blind spot" in the retina and an equivalent blind spot in our field of vision.

The intricate structure of the retina comprises multiple layers populated with different cells, all of which play a critical role in the initial interpretation of visual cues. The photoreceptors, specifically rods and cones, are sensitive to light energy, and this sensitivity prompts a shift in their membrane potential. This shift subsequently determines the amount of neurotransmitter released onto the bipolar cells in the outer synaptic layer. The bipolar cell is the intermediary between a photoreceptor and a retinal ganglion cell (RGC) in the inner synaptic layer within the retina. The processing within the retina is assisted by Amacrine cells before the RGC generates an action potential. The axons of RGCs, nestled in the innermost layer of the retina, converge at the optic disc, exiting the eye as the optic nerve. Due to the course these axons take through the retina, the rear of the eye, where the optic nerve originates, is devoid of photoreceptors. This results in a "blind spot" in the retina, reflecting an identical blind spot in our visual field.

It is important to note that the photoreceptors (rods and cones) within the retina are behind the axons, RGCs, bipolar cells, and retinal blood vessels. These structures absorb a considerable amount of light before it reaches the photoreceptor cells. Yet, the fovea is at the retina's center — a small area devoid of supporting cells and blood vessels, only housing photoreceptors. As such, visual acuity — the clarity of vision — is optimal at the fovea due to minimal absorption of incoming light by other retinal structures. As one moves away from the foveal center in any direction, there is a noticeable drop in visual acuity. Each of the fovea's photoreceptor cells is connected to a single RGC. It follows that the RGC does not need to amalgamate inputs from multiple photoreceptors, enhancing the precision of visual transduction.

Conversely, at the peripheries of the retina, several photoreceptors converge on RGCs (via the bipolar cells) in ratios as high as 50 to 1. The disparity in visual acuity between the fovea and the peripheral retina is starkly evident — focus on a word positioned in the middle of this paragraph without moving your eyes, and words at the beginning or end appear blurry and out of focus. The peripheral retina is responsible for creating the images in your peripheral field of view; however, these images often have indistinct, fuzzy edges, and the words must be more clearly discernible. So, a significant portion of the neural function of the eyes is concentrated on moving the eyes and head to ensure important visual stimuli are centered on the fovea.

Photoreceptors, the cells responsible for capturing light in the eye, are composed of two distinct components: the internal and external segments. The former harbors the nucleus and various other cell organelles, while the latter is a niche area enabling photoreception. Two distinct photoreceptor types exist — rods and cones — characterized by the morphology of their external segments. The rods — named for their rod-like segments — house membranous disks filled with the light-sensitive pigment rhodopsin. The cone photoreceptors, on the other hand, hold their light-sensitive pigments within the cell membrane's invaginations, and their external segments take on a conical shape. Cone photoreceptors possess three photopigments, namely opsins, each responsive to a specific light wavelength. The color of visible light is determined by its wavelength, and the photopigments in the human eye are adept at discerning three fundamental colors: red, green, and blue.

 Core: Anatomy and Physiology

Regulation of Hormone Secretion

JoVE 14978

Regulation of hormone secretion is a finely tuned orchestration driven by various types of stimuli, encompassing neural, humoral, and hormonal signals. Environmental cues instigate neural stimuli, where action potentials traverse nerve fibers to reach their designated targets. An illustrative scenario is the body's response to stress, wherein the sympathetic nervous system releases epinephrine from the adrenal glands, inducing the well-known 'fight or flight' reaction.

Humoral stimuli, conversely, involve concentration fluctuations of specific ions or nutrients in the bloodstream. A case in point is the response to diminished blood calcium levels, triggering the secretion of parathyroid hormones to bolster circulating calcium concentrations.

Hormonal stimuli manifest when one hormone induces the secretion of another hormone from a distinct endocrine organ. For instance, the hypothalamus secretes the thyrotropin-releasing hormone, instigating the anterior pituitary lobe to release the thyroid-stimulating hormone (TSH). TSH, in turn, stimulates the thyroid glands to discharge thyroid hormones like T3. As circulating T3 levels rise, a sophisticated negative feedback loop comes into play. This loop inhibits the continuous production of T3 by dispatching inhibitory signals to both the pituitary gland and hypothalamus. The endocrine system maintains a delicate balance in this intricate interplay of signals and responses, ensuring precise regulation and harmony in hormone secretion.

 Core: Anatomy and Physiology

Nervous Tissue: Glial Cells

JoVE 14996

Glia, or neuroglia, are vital support cells that assist neurons in their functions. The term "glia" originates from the Greek word for "glue," reflecting their role in holding the nervous system together. These cells can be categorized into six types: four in the central nervous system (CNS) and two in the peripheral nervous system (PNS).

The CNS glial cell includes the astrocytes, the oligodendrocytes, the microglia, and the ependymal cells.

Astrocytes are star-shaped glial cells that interact with neurons, blood vessels, and the pia mater, a connective tissue covering the CNS. Their primary role is maintaining chemical balance, removing excess signaling molecules, contributing to the blood-brain barrier, and guiding neuronal growth and connections during embryonic development. Oligodendrocytes, smaller than astrocytes, form and maintain the myelin sheath around CNS axons. This myelin sheath is a multilayered lipid and protein covering that insulates axons. Microglial cells or microglia function as phagocytes, removing cellular debris and damaged nervous tissue. Ependymal cells, arranged in a single layer, line the brain's ventricles and the spinal cord's central canal. These cells produce, monitor, and help circulate cerebrospinal fluid. They also form the blood-cerebrospinal fluid barrier.

In the PNS, Schwann cells encircle axons, similar to oligodendrocytes, forming the myelin sheath around them. However, while a single oligodendrocyte can myelinate several axons, each Schwann cell myelinates only one axon. Finally, satellite cells surround the cell bodies of neurons in the PNS ganglia. They provide structural support and regulate material exchange between neuronal cell bodies and interstitial fluid in the PNS.

 Core: Anatomy and Physiology

Charge and Current

JoVE 15063

Electric charge is the most fundamental quantity in an electric circuit. The effects of electric charge are encountered daily, such as when a wool sweater sticks to the human body or when a person receives a shock while walking on a carpet.

Charge is an inherent property of the atomic particles that make up matter and is measured in units called coulombs (C). Matter is composed of atoms, each consisting of electrons, protons, and neutrons. Electrons have a negative charge (-e), while protons carry a positive charge (+e), with both charges having the same magnitude. An atom is electrically neutral when it has an equal number of protons and electrons.

Key points about electric charge include:

  • • The coulomb (C) is a relatively large unit for measuring charges, and practical charge values are often in picocoulombs (pC), nanocoulombs (nC), or microcoulombs (μC).
  • • Observations reveal that natural charges exist only as integral multiples of the elementary charge (e), the charge carried by an electron or proton.
  • • The law of conservation of charge asserts that electric charge cannot be created nor destroyed; it can only be transferred. As a result, the total electric charge in a system remains constant. Charge can be moved from one place to another and converted into different forms of energy.

The movement of positive charges in one direction and negative charges in the opposite direction gives rise to electric current. Typically, a conventional current flow is defined as the movement of positive charges, and it is measured in units known as amperes (A). When the current remains constant over time, it is called direct current (DC). In cases where the current varies with time, it is denoted as a time-varying current, often referred to as alternating current (AC). AC is commonly used in households to power appliances like air conditioners, refrigerators, and washing machines.

 Core: Electrical Engineering

Sum and Difference OpAmps

JoVE 15079

Operational amplifiers (op-amps) are versatile devices that extend beyond amplification. In this context, two specific op-amp configurations are explored: the summing and difference amplifiers.

A summing amplifier, or an adder, utilizes an op-amp to merge multiple input signals into a single output signal. When audio signals are introduced into its input channels, the input resistors initiate currents that traverse feedback resistors, resulting in an output voltage. Applying Kirchhoff's current law and Ohm's Law, the output voltage becomes a weighted sum of the input signals. This basic differential op-amp circuit plays a role in mitigating audio track noise in audio mixer applications.

Equation1

Conversely, a difference amplifier, often called a subtractor, amplifies the disparity between two input signals. Two equations are established by applying Kirchhoff's current law at the inverting and non-inverting nodes. An output voltage expression is derived when equal voltages exist across both nodes. However, specific resistance ratio conditions must be satisfied for the difference amplifier to nullify common signals between the inputs. When the resistances in the corresponding branches are equal, the difference amplifier functions as a subtractor.

Equation2

In summary, operational amplifiers offer amplification and the ability to perform addition and subtraction tasks. The summing amplifier combines inputs into a weighted sum. The difference amplifier amplifies the disparity between two inputs and rejects common signals. Understanding these configurations opens the door to many practical applications in signal processing and circuit design.

 Core: Electrical Engineering

RL Circuit with Source

JoVE 15096

When an RL (Resistor-Inductor) circuit is connected to a DC source, the complete response of the circuit can be divided into two parts: the transient response and the steady-state response.

The transient response of the circuit is its temporary reaction to the sudden application of the DC source. This response is characterized by a current that exponentially decays to zero as time approaches infinity. During this transitional period, the inductor behaves like a short circuit, causing the source voltage to drop across the resistor.

Once the transient response has decayed to zero, the circuit enters its steady-state phase. In this phase, the current in the circuit becomes stable, equaling the ratio of the source voltage to the resistance of the circuit. This is the steady-state response.

The constant term in the transient response can be determined by substituting the initial current through the inductor at the time the switch is closed (t=0). This gives us the initial condition for the transient response.

Plotting the complete step response of the circuit shows that the initial current decreases exponentially until it reaches a steady-state value. From this current response, the voltage response is derived, which also follows an exponential decay from a maximum value to zero.

If the initial current in the circuit is zero, the complete step response shows the current increasing exponentially until it reaches the steady-state value. Correspondingly, the voltage response starts at the source voltage and decreases exponentially to zero.

In conclusion, understanding the complete response of an RL circuit to a DC source provides valuable insights into how these circuits react to sudden changes in input voltage. This knowledge is essential for designing and analyzing electronic circuits, particularly in applications such as power supply filtering and signal processing, where inductors are used extensively.

 Core: Electrical Engineering

Thévenin Equivalent Circuits

JoVE 15113

The household power distribution system, encompassing distribution lines and transformers, serves as the primary network. Electrical appliances within a household can be represented as load impedance. To simplify this intricate distribution system, Thévenin's theorem can be applied to create a Thévenin equivalent circuit. If an AC circuit is partitioned into two parts (circuit A and circuit B), connected by a single pair of terminals as shown in Figure 1.

Figure1

Figure 1: Circuit portioned into two parts

Figure2

Figure 2:Circuit A replaced by its Thévenin equivalent circuit

Replacing circuit A with its Thévenin equivalent circuit (a voltage source in series with an impedance) does not alter the current or voltage of any element in circuit B (shown in Figure 2). The values of the currents and voltages of all the circuit elements in circuit B will be the same irrespective of whether circuit B is connected to circuit A or its Thévenin equivalent. Two parameters are required to find the Thévenin equivalent circuit: the Thévenin voltage and the Thévenin impedance. Figure 3 shows an open circuit connected across the terminals of circuit A to determine the open-circuit voltage Voc , while Figure 4 indicates that the Thévenin impedance Zt is the equivalent impedance of circuit A*.

Figure3

Figure 3:Thévenin equivalent circuit with Voc .

Figure4

Figure 4:Thévenin equivalent circuit showing Zt. .

Circuit A* is formed from circuit A by replacing all independent voltage sources with short circuits and all independent current sources with open circuits. Generally, the Thévenin impedance Zt can be determined by replacing series or parallel impedances with equivalent impedances repeatedly.

 Core: Electrical Engineering

Titration of a Weak Acid with a Strong Base

JoVE 17361

In titrating a weak acid with a strong base, different calculation methods are applied at various stages. Initially, the pH of a weak acid like acetic acid is calculated using its dissociation constant (Ka) and an ICE table. Upon addition of a strong base such as sodium hydroxide, a buffer forms, and its pH is determined using the Henderson-Hasselbalch equation. As more base is added and the titration reaches the halfway point, the pH becomes equal to the pKa of the acid, indicating equal concentrations of the acid and its conjugate base. At the equivalence point, all the acid is converted to its conjugate base, and the pH is calculated using the base's dissociation constant (Kb) and an ICE table. Beyond the equivalence point, the pH is governed by the concentration of the excess strong base.

 Core: Analytical Chemistry

Yeast Maintenance

JoVE 5095

Research performed in the yeast Saccharomyces cerevisiae has significantly improved our understanding of important cellular phenomona such as regulation of the cell cycle, aging, and cell death. The many benefits of working with S. cerevisiae include the facts that they are inexpensive to grow in the lab and that many ready-to-use strains are now commercially available. Nevertheless, proper maintenance of this organism is critical for successful experiments.

This video will provide an overview of how to grow and maintain S. cerevisiae in the lab. Basic concepts required for monitoring the proliferation of a yeast population, such as how to generate a growth curve using a spectrophotometer, are explained. This video also demonstrates the hands-on techniques required to maintain S. cerevisiae in the lab, including preparation of media, how to start a new culture of yeast cells, and how to store those cultures. Finally, the video shows off some of the ways these handling and maintenance techniques are applied in scientific research.

 Biology I

Chick ex ovo Culture

JoVE 5157

One strength of the chicken (Gallus gallus domesticus) as a model organism for developmental biology is that the embryo develops outside the female and is easily accessible for experimental manipulation. Many techniques allow scientists to examine chicken embryos inside the eggshell (in ovo), but embryonic access can be limited at later stages of development. Fortunately, chicks can also be cultured ex ovo, or outside of the eggshell. The major advantage to ex ovo culture is greater access to tissues that might otherwise be obstructed by the shell or the orientation of the chick within the egg, especially for embryos in later stages of development.

There are two principle strategies to ex ovo culture: whole yolk culture and explant culture. During whole yolk culture, the eggshell is cracked and the contents are transferred to a simple housing vessel. However, in explant culture methods, the embryo is excised from the yolk and mounted in the housing vessel to maintain membrane tension, which is important for normal development.

Basic protocols for whole-yolk and explant techniques will be provided in this video, along with a discussion of the pros and cons of culturing chicks outside of the shell. Finally, experimental applications of ex ovo culture will be discussed, demonstrating how this approach is used to improve access to the embryo for microscopy and genetic manipulation of late stage embryos.

 Biology II

An Introduction to Behavioral Neuroscience

JoVE 5210

Behavioral neuroscience is the study of how the nervous system guides behavior, and how the various functional areas and networks within the brain correlate to specific behaviors and disease states. Researchers in this field utilize a wide variety of experimental methods ranging from complex animal training techniques to sophisticated imaging experiments in human subjects.

This video first offers a historical overview of some of the major milestones that lead to our current understanding of the brain’s control over behavior. Then, some of the fundamental questions asked by behavioral neuroscientists are presented, which all involve the study of neural correlates, or specific brain regions whose activation is responsible for a given function. Next, prominent methods used to answer those questions are reviewed for both human and animal subjects, such as operant conditioning and functional neuroimaging. Finally, experimental applications of these techniques are presented, including animal training using a Skinner box, and the use of electroencephalography to investigate human neurological disease.

 Neuroscience

Embryonic Stem Cell Culture and Differentiation

JoVE 5332

Culturing embryonic stem (ES) cells requires conditions that maintain these cells in an undifferentiated state to preserve their capacity for self-renewal and pluripotency. Stem cell biologists are continuously optimizing methods to improve the efficiency of ES cell culture, and are simultaneously trying to direct the differentiation of ES cells into specific cell types that could be used in regenerative medicine.

This video describes the basic principles of ES cell culture, and demonstrates a general protocol to grow and passage ES cells. We also take a closer look at the hanging drop method, which is used to differentiate ES cells. Finally, this video will describe a few applications of ES cell culture and differentiation techniques, including a method used to generate functional heart muscle cells in vitro.

 Developmental Biology

An Introduction to Motor Control

JoVE 5422

Motor control involves integration and processing of sensory information by our nervous system, followed by a response through our skeletal system to perform a voluntary or involuntary action. It is vital to understand how our neuroskeletal system controls motor behavior in order to evaluate injuries pertaining to general movement, reflexes, and coordination. An improved understanding of motor control will help behavioral neuroscientists in developing useful tools to treat motor disorders, such as Parkinson's or Huntington's disease.

This video briefly reviews the neuroanatomical structures and connections that play a major role in controlling motion. Fundamental questions currently being asked in the field of motor control are introduced, followed by some of the methods being employed to answer those questions. Lastly, the application sections reviews a few specific experiments conducted in neuroscience labs interested in studying this phenomenon.

 Behavioral Science

Genome Editing

JoVE 5554

A well-established technique for modifying specific sequences in the genome is gene targeting by homologous recombination, but this method can be laborious and only works in certain organisms. Recent advances have led to the development of “genome editing”, which works by inducing double-strand breaks in DNA using engineered nuclease enzymes guided to target genomic sites by either proteins or RNAs that recognize specific sequences. When a cell attempts to repair this damage, mutations can be introduced into the targeted DNA region.

In this video, JoVE explains the principles behind genome editing, emphasizing how this technique relates to DNA repair mechanisms. Then, three major genome editing methods—zinc finger nucleases, TALENs, and the CRISPR-Cas9 system—are reviewed, followed by a protocol for using CRISPR to create targeted genetic changes in mammalian cells. Finally, we discuss some current research that applies genome editing to alter the genetic material in model organisms or cultured cells.

 Genetics

An Introduction to Cell Death

JoVE 5649

Necrosis, apoptosis, and autophagic cell death are all manners in which cells can die, and these mechanisms can be induced by different stimuli, such as cell injury, low nutrient levels, or signaling proteins. Whereas necrosis is considered to be an “accidental” or unexpected form of cell death, evidence exists that apoptosis and autophagy are both programmed and “planned” by cells.

In this introductory video, JoVE highlights key discoveries pertaining to cell death, including recent work done in worms that helped identify genes involved in apoptosis. We then explore questions asked by scientists studying cell death, some of which look at different death pathways and their interactions. Finally, several methods to assess cell death are discussed, and we note how researchers are applying these techniques in their experiments today.

 Cell Biology

Enzyme Assays and Kinetics

JoVE 5692

Enzyme kinetics describes the catalytic effects of enzymes, which are biomolecules that facilitate chemical reactions necessary for living organisms. Enzymes act on molecules, referred to as substrates, to form products. Enzyme kinetic parameters are determined via assays that directly or indirectly measure changes in substrate or product concentration over time. 

This video will cover the basic principles of enzyme kinetics (including rate equations) and kinetic models. The concepts governing enzyme assays are also discussed, followed by a typical colorimetric assay. The applications section discusses an enzyme assay via Förster resonance energy transfer (FRET) analysis, characterizing extracellular enzyme activity in the environment, and investigating DNA repair kinetics using molecular probes.

Enzymes are biochemical catalysts that are essential for life. Enzyme assays are used to study the kinetic properties of enzymatic reactions, elucidating the catalytic effects of enzymes. This video will cover enzyme kinetics and assays, go over a general procedure, and show some applications.

Enzymes are proteins, or less often RNAs, that act on a specific reactant, referred to as the substrate. An enzyme reduces the activation energy needed to initiate a biochemical reaction, causing the reaction to occur at a faster rate.

Enzymatic reactions can be broken up into three elementary components. The first is the formation of the enzyme-substrate complex, formed by the binding of the substrate to the enzyme active site. The complex can decompose into its original constituents. This is the second elementary reaction. Alternatively, the complex can form the product and recover the enzyme, the third elementary reaction.

The kinetics of an elementary reaction is given by the elementary rate law equation. Rate law equations give the rate in terms of the concentration of the reactants and a rate constant. Each of the elementary reactions has an individual rate law equation, with its own rate constant. These equations can be distilled down to a kinetic model known as the Michaelis-Menten equation. This gives the reaction rate in terms of the substrate concentration; which can be experimentally determined. Some general trends for enzyme reactions can be identified using the Michaelis-Menten equation. At high substrate concentration, a saturation point is reached, called Vmax. Here, the rate is limited by the total enzyme concentration, and the number of substrate molecules an enzyme converts into product per given time, also known as kcat. In Michaelis-Menten kinetics kcat is one of the two constants that govern reaction rate. The other constant, KM, is known as the affinity constant. KM is also equivalent to the concentration where the reaction rate is equivalent to one-half Vmax . An enzyme with a higher affinity will have a lower KM and reach Vmax faster, while an enzyme with lower affinity will have a higher KM and take longer to reach Vmax. Knowing kcat and KM allows for enzymes to be compared. To do this we use a ratio called enzyme efficiency. Higher kcat and lower KM result in higher efficiencies, while lower kcat and higher KM results in lower.

The factors used to elucidate enzyme kinetics must be determined experimentally. These assays are typically performed by mixing an enzyme and substrate solution in a controlled environment. Observations are made by measuring the changes in concentration of the substrate, product, or byproducts with respect to time.

The change in concentration over time is used to determine the reaction rate. In order to determine the kinetics, rate data must be obtained at multiple concentrations. If a plot of the inverse initial rate vs. inverse initial concentration, known as the Lineweaver-Burk plot, is linear, then the reaction follows Michaelis-Menten kinetics. The slope and intercept of the line allow for the determination of the kinetic parameters KM and Vmax, which can then be used to calculate kcat and the enzyme efficiency.

Now that the principles of enzyme kinetics have been discussed, let's look at how a typical enzyme assay is performed.

In this procedure a colorimetric assay is demonstrated. The first step is to generate a standard curve, which will correlate absorbance with substrate concentration. Solutions of known concentration are prepared along with a control sample. A developer solution that reacts with the substrate is added to produce a colored compound. Absorbance is measured and plotted against concentration to generate the standard curve.

To perform the assay, a known concentration of substrate is prepared along with the appropriate amount of enzyme. The enzyme and substrate are mixed and allowed to incubate for a set time interval. pH and temperature are controlled with buffer solutions and heating blocks. A quenching agent is added to stop the reaction. Developer solution is then added to the reactions and mixed. The solutions are then placed in cuvettes and absorbance is measured. The amount of substrate consumed is determined by comparing the measured absorbance to the standard curve. Using the collected data, initial reaction rates are determined by plotting concentration over time. Finally, with the rate data and concentration, the Michaelis-Menten plot is made. This allows for the determination of kinetic properties for the enzyme such as turnover number and enzyme efficiency.

Now that we've reviewed an assay procedure, let's look at other ways assays are performed and their applications.

In this procedure FRET analysis is used to study the kinetics of a protease hydrolyzing a peptide bond of a protein. These emissions can be measured, allowing for a continuous and quantitative analysis of substrate consumption and production, aiding in the determination of the reaction kinetics.

Enzyme assays can be used in environmental science to determine the levels of extracellular enzyme activity in the environment. Waters, soils, and sediments can be collected from the environment and processed in the laboratory. Extracellular enzymatic activity of these materials can then be characterized using enzyme assays. This is a useful tool for understanding how the environment processes organic material.

A cell's DNA repair mechanism can be evaluated by studying the kinetics of enzymes found in the nucleus. The rate at which an enzyme removes DNA lesions, or damages, can be measured using fluorescent molecular beacons, which only fluoresce when bound to unique DNA sequences. The level of DNA repair can be measured in real time by detecting the fluorescently labeled cleavage products.

You've just watched JoVE's video on enzyme kinetics and assays. This video explained enzyme kinetics, covered assay concepts, went over a general procedure, and described some applications.

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 Biochemistry

Overview of Bioprocess Engineering

JoVE 5791

Bioprocessing is a method that uses living organisms to produce a desired target product. Often, bioprocessing refers to the use of bioreactors to produce protein products from genetically engineered organisms. This field is responsible for the large-scale manufacture of biotherapeutics; drugs that have become essential to improving the quality of life for many with complex diseases like cancer, autoimmune diseases and HIV/AIDS.

This video will introduce the engineering approach to designing a targeted protein-production system. The prominent methods in the field, as well as some key challenges, and applications of the technology are also considered.

 Bioengineering

Bond Dissociation Energy and Activation Energy

JoVE 11700

Bond energy is the energy required to break a bond homolytically. These values are usually expressed in units of kcal/mol or kJ/mol and are referred to as bond dissociation energies when given for specific bonds or average bond energies when indicated for a given type of bond over many compounds. Firstly, the bond dissociation energy for a single bond is weaker than that of a double bond, which in turn is weaker than that of a triple bond. Secondly, hydrogen forms relatively strong bonds with carbon, nitrogen, and oxygen. Finally, with the exception of carbon and hydrogen, single bonds between atoms of the same element are relatively weak. Reactions between organic compounds involve the making and breaking of bonds. Hence, the strengths of bonds and their resistance to breaking are essential concepts in organic chemistry.

Reactions in which bonds are broken pass through a high-energy transition state before transforming into products. In order to reach this transition state, the reactant molecules must be oriented in a suitable direction and must be supplied with certain threshold energy. The activation energy, ΔG, is the energy provided to the reactants to raise them to the transition state. Overall, for a reaction to occur, the reacting molecules must collide or otherwise interact. The necessary activation energy for the reactant–transition-state jump is provided by the kinetic energy of the colliding particles. Lastly, the colliding molecules must collide in a specific orientation so as to maximize the impact of the collision.

 Core: Organic Chemistry

Ethers to Alkyl Halides: Acidic Cleavage

JoVE 11741

Ethers are generally unreactive and unsuitable for direct nucleophilic substitution reactions since the alkoxy groups are strong bases and, therefore, poor leaving groups. However, ethers readily undergo acidic-cleavage reactions. Ethers can be converted to alkyl halides when heated with strong acids such as HBr and HI in a sequence of two substitution reactions.

Figure1

In the first step, the ether is converted into an alkyl halide and alcohol.

Figure2

In the second step, the alcohol reacts with the excess HX acid to form another equivalent of an alkyl halide.

Figure3

Acidic cleavage of ethers can occur either by the SN1 mechanism or SN2 mechanism, depending on the substrate. Ethers with primary and secondary alkyl groups as substrates react with the SN2 mechanism in which the nucleophile attacks the protonated ether at the less hindered site. Ethers with a tertiary, benzylic, or allylic group undergo the SN1 mechanism because the substrates can produce more stable intermediate carbocations. Reactivity of the halogen acids with ethers increases relative to the nucleophilicity of the halide ions. HI and HBr are more reactive, HCl is less efficient, and HF is unsuitable.

 Core: Organic Chemistry

Oxymercuration-Reduction of Alkenes

JoVE 11778

Oxymercuration–reduction of alkenes is one of the major reactions converting alkenes to alcohols. It involves the hydration of alkenes with mercuric acetate in a mixture of tetrahydrofuran and water, forming an organomercury adduct. This is followed by a demercuration step in which the adduct is reduced to an alcohol using sodium borohydride.

Figure1

In the mixture of water and tetrahydrofuran, tetrahydrofuran acts as a solvent dissolving the alkene and the aqueous mercuric acetate solution, while water functions as a reactant and a solvent for mercuric acetate.

Oxymercuration–Demercuration Mechanism

Consider the conversion of 2-methyl-2-butene to yield 2-methyl-2-butanol.

The mechanism proceeds with the dissociation of mercuric acetate, forming an electrophilic mercuric cation and an acetate anion. The alkene π bond attacks the electrophilic mercuric cation, resulting in a bridged-mercurinium-ion intermediate.

Figure2

Regiochemical and Stereochemical Outcome

The bridged-mercurinium-ion intermediate is a resonance hybrid of a carbocation and a bridged mercurinium ion. The partial positive charge is shared between the more substituted carbon atom and the mercury atom, minimizing the chance of a carbocation rearrangement. Furthermore, the carbon–mercury bond to the more substituted carbon is longer and can be easily broken.

Figure3

The factors mentioned above lead to the nucleophilic attack by water exclusively at the more substituted carbon, opening the three-membered ring.

The oxymercuration step is stereospecific, as the attack by water on the bridged mercurinium ion leads to the anti addition of the hydroxyl group. A proton transfer completes the oxymercuration step, forming an organomercury compound.

Lastly, the oxymercuration adduct is treated with sodium borohydride through a process called demercuration to yield an alcohol with Markovnikov's orientation.

During the demercuration step, as the hydrogen can replace the mercury species in either a syn or an anti fashion with respect to the hydroxyl group, the overall reaction produces a racemic mixture of two enantiomeric alcohols.

 Core: Organic Chemistry

Electrophilic Addition to Alkynes: Halogenation

JoVE 11836

Introduction

Halogenation is another class of electrophilic addition reactions where a halogen molecule gets added across a π bond. In alkynes, the presence of two π bonds allows for the addition of two equivalents of halogens (bromine or chlorine). The addition of the first halogen molecule forms a trans-dihaloalkene as the major product and the cis isomer as the minor product. Subsequent addition of the second equivalent yields the tetrahalide.

Figure1

Reaction Mechanism

In the first step, a π bond from the alkyne acts as a nucleophile and attacks the electrophilic center on the polarized halogen molecule, displacing the halide ion and forming a cyclic halonium ion intermediate. In the next step, a nucleophilic attack by the halide ion opens the ring and forms the trans-dihaloalkene. Since the nucleophile attacks the halonium ion from the backside, the net result is an anti addition where the two halogen atoms are trans to each other.

Figure2

The addition of a second equivalent of halogen across the alkene π bond also proceeds via the formation of a bridged halonium ion to give the tetrahalide as the final product.

Figure3

For example, the addition of bromine to 2-butyne in the presence of acetic acid and lithium bromide favors anti addition and preferentially forms the trans or (E)-2,3-dibromo-2-butene as the major product. The corresponding cis isomer, (Z)-2,3-dibromo-2-butene, is formed in lower yields. A second addition gives 2,2,3,3-tetrabromobutane.

Figure4

Reactivity of alkynes and alkenes towards electrophilic addition

Alkynes are less reactive than alkenes towards electrophilic addition reactions. The reasons are twofold. First, the carbon atoms of a triple bond are sp hybridized in contrast to the double bonds that are sp2 hybridized. Since the sp hybrid orbitals have a higher s-character and are more electronegative, the π electrons in C≡C are held more tightly than in C=C. As a result, in alkynes, the π electrons are not readily available for the nucleophilic attack, making them less reactive towards electrophilic addition than alkenes.     

Secondly, the cyclic halonium ion formed from alkynes is a three-membered ring with a double bond where the 120° bond angle of an sp2 carbon is constrained into a triangle.

Figure5a Figure5b
Alkyne halonium ion Alkene halonium ion

In contrast, the cyclic intermediate in alkenes is a three-membered ring with an sp3 hybridized carbon where a bond angle of 109° is constrained into a triangle. Therefore, the larger ring strain associated with the alkyne halonium ions makes them more unstable and hinders their formation. 

 Core: Organic Chemistry

Drugs that Stabilize Microtubules

JoVE 11915

Microtubules are dynamic structures that undergo cycles of catastrophe and rescue. The microtubules play a central role in cell division by forming the spindle apparatus for segregating the chromosomes. This makes them ideal targets for regulating dividing cells in tumors and malignant cancer cells. Microtubule stabilizing drugs help stabilize the microtubule formation and promote its polymerization. Paclitaxel was the first microtubule stabilizing agent used as anticancer drug in chemotherapy after medical approval in 1993.

Some important microtubule-stabilizing drugs include taxanes, epothilones, laulimalide, and dictyostatin. These drugs can be used to treat tumors in breast, lung, prostate, and ovarian cancers. Taxanes are the widely used class of microtubule-stabilizing drugs that target microtubules in the spindle apparatus and cytoskeleton and disrupt cellular processes, including cell division. If given in low concentration, these drugs lead to the formation of multipolar spindle apparatus during the G1 phase, a process called mitotic slippage. This leads to cell cycle arrest and apoptosis. If the drugs are given in a higher dosage, it causes mitotic arrest of cells during the G2/M phase of mitosis. The mitotic arrest prevents chromosome segregation, resulting in tetraploid G1 cells, with eventual apoptosis of cancer or tumor cells. However, these drugs have certain limitations because ABC transporters overexpressed in tumor, which pump them out of the cell.

These drugs have specific binding sites on the tubulin dimers. Taxane sites are hydrophobic pockets between the lateral interface of two adjacent protofilaments in the microtubule lumen. Epothilones used in paclitaxel-resistant tumors also use hydrophobic taxane sites for binding with microtubules.

 Core: Cell Biology

Base-Catalyzed Ring-Opening of Epoxides

JoVE 12110

Due to their highly strained structures, epoxides can readily undergo ring-opening reactions through nucleophilic substitution, either in the presence of an acid or a base. The nucleophilic substitution reactions in the presence of acid are called acid-catalyzed ring-opening reactions, and nucleophilic substitution reactions in the presence of a base are called base-catalyzed ring-opening reactions. Epoxides undergo base-catalyzed ring-opening reactions in the presence of a strong nucleophile or a base. A variety of nucleophiles like sodium hydroxide, sodium alkoxide, sodium hydrosulfide, sodium cyanide, lithium aluminum hydride, and Grignard reagent can open the epoxide ring. The nucleophilic attack on the epoxide ring proceeds via an SN2 mechanism and involves an alkoxide intermediate. The product formed in base-catalyzed reactions shows SN2-like stereoselectivity and regioselectivity. The nucleophile attacks at the less-hindered carbon and anti to the leaving group, leading to the inversion of configuration at a chiral center. In contrast to epoxides, acyclic ethers do not undergo direct nucleophilic substitutions, as the reaction is thermodynamically unfavorable.

 Core: Organic Chemistry

[4+2] Cycloaddition of Conjugated Dienes: Diels–Alder Reaction

JoVE 12417

The Diels–Alder reaction is an example of a thermal pericyclic reaction between a conjugated diene and an alkene or alkyne, commonly referred to as a dienophile. The reaction involves a concerted movement of six π electrons, four from the diene and two from the dienophile, forming an unsaturated six-membered ring. As a result, these reactions are classified as [4+2] cycloadditions.

Figure1

From a molecular orbital perspective, the interacting lobes of the two π systems must be in phase to permit the formation of new σ bonds in a synchronous manner. For molecules in the ground state, the interaction between HOMO (diene) and LUMO (dienophile) or HOMO (dienophile) and LUMO (diene) satisfies the orbital symmetry requirements. The two π components interact suprafacially in each case, making Diels–Alder a thermally allowed [4+2] cycloaddition reaction.

Figure2

 Core: Organic Chemistry

Multipotency of Hematopoietic Stem Cells

JoVE 12512

The hematopoietic stem cells or HSCs are multipotent, meaning they can differentiate and give rise to all blood and immune cells. HSCs are maintained in the quiescent stage until an external stimulus initiates their differentiation. The multipotent HSCs exist as two heterogeneous populations, long-term repopulating cells (LTRC) and short-term repopulating cells (STRC). The two HSC populations have different surface markers or receptors and are classified based on quiescence and long-term renewal capacity.

The long-term HSCs rarely divide as they maintain long gaps between successive cell divisions while keeping their metabolic activity to the bare minimum. Thus, the LTRCs undergo a few rounds of symmetric cell divisions before entering the state of dormancy. This allows them to expand the number of HSCs and maintain a stem cell pool before exhausting their regenerative and self-renewal potential.

LTRCs produce the more active multipotent short-term repopulating cells (STRC) in response to external stimuli. STRCs undergo fewer self-renewal divisions than the LTRCs and differentiate into specific blood or immune cells. LTRCs and STRCs find application in engraftment and therapeutic studies. The LTRCs can sustain hematopoiesis for upto four months upon transplantation into the recipients, while transplantation of STRCs supports hematopoiesis for only a few weeks. A delicate balance between dormancy and differentiation stimuli maintains the HSC population for the long run. A slow or non-responsive quiescent HSCs leave the body devoid of differentiated blood cells such as erythrocytes, phagocytes, and lymphocytes which are necessary for transporting nutrients or providing immune surveillance. In contrast, highly active HSCs exhaust the stem cell population making them unavailable to renew, repair, or replace blood cells lost due to an injury or infection.

 Core: Cell Biology

Alkylation of β-Diester Enolates: Malonic Ester Synthesis

JoVE 13075

Malonic ester synthesis is a method to obtain α substituted carboxylic acids from ꞵ-diesters such as diethyl malonate and alkyl halides.

Figure1

The reaction proceeds via abstraction of the acidic α hydrogen from a ꞵ-diester to produce a doubly stabilized enolate ion. The nucleophilic enolate attacks the alkyl halide in an SN2 manner to form an alkylated malonic ester intermediate with a new C–C bond. Further treating the intermediate with aqueous acid or base results in the hydrolysis of the two ester groups to give a 1,3-dicarboxylic acid. The resulting ꞵ-diacid is unstable at high temperatures and readily eliminates CO2 through a cyclic six-membered transition state, forming an enol. The enol tautomerizes to its more stable keto form producing a monosubstituted carboxylic acid. However, a disubstituted carboxylic acid is achieved if the deprotonation and alkylation steps are repeated before hydrolysis and decarboxylation.

 Core: Organic Chemistry

Principles Of Column Chromatography

JoVE 13376

The chromatography technique was first invented in 1901 by Michael S. Tswett, a Russian botanist, to separate plant pigments using organic solvents. Further, in 1941, Archer John Porter Martin and R. L. M. Synge modified the technique by packing silica gel into a column. A mixture of amino acids was then separated on the packed column using chloroform and water mixture as the mobile phase. This was the first report on column chromatography. At present, column chromatography is a widely used technique for separating various types of compounds from a sample mixture.

Factors Influencing Efficient Separation of Proteins

Various parameters like column material, packing, and operational conditions such as flow rates and temperature determine the efficiency of separation by column chromatography.

The choice of column material or matrix determines the extent of interaction with the sample. The matrix material must be tightly and uniformly packed in the column. Air bubbles, debris, large particles, and precipitates interfere with the uniform flow of the solvent through the column, affecting its separation efficiency. The column should also be free from particulate matter.

The sample injected into the column should be clear and free from aggregates that might clog the column, hindering solvent flow. The flow rate of solvent also affects the separation. Very high or very low flow rates of solvents result in inefficient separation of compounds and impure preparations. Very high rates may also disturb the column packing, affecting process efficiency. In addition, the composition of elution buffer is also an important factor. It should be non-corrosive and compatible with the sample and also the column material to prevent in situ precipitation or dissolution.

Another operational parameter, temperature, also plays an important role in the process. It decides the stability of the sample, column material, and solvent buffer. Also, a constant temperature throughout the column efficiently resolves compounds. After completing the separation process, the columns must be washed thoroughly by repeatedly passing a suitable solvent to avoid sample contamination in further runs. Occasionally, the solvent is passed in a column in the reverse direction to remove any clogged material.

Limitations

Though a very widely used technique, the method still has some limitations. It is a very time-consuming method as the flow rates need to be slower for better resolution of the compounds. Also, large quantities of highly pure solvents required in the mobile phase make the process expensive. This also raises the scaling-up cost when higher yields of pure compounds are needed.

 Core: Cell Biology

Immunofluorescence Microscopy

JoVE 13392

A fluorescence microscope uses fluorescent chromophores called fluorochromes, which can absorb energy from a light source and then emit this energy as visible light. Fluorochromes include naturally fluorescent substances (such as chlorophylls) and fluorescent stains that are added to the specimen to create contrast. Dyes such as Texas red and FITC are examples of fluorochromes. Other examples include the nucleic acid dyes 4’,6’-diamidino-2-phenylindole (DAPI), and acridine orange.

The microscope irradiates the sample with short wavelength excitation, such as ultraviolet or blue light. The chromophores absorb the excitation light and emit visible light of longer wavelengths. The excitation light is then filtered out (partly because ultraviolet light is harmful to the eyes) so that only visible light passes through the ocular lens, producing an image of the specimen in bright colors against a dark background.

Fluorescence microscopes can identify pathogens, find particular species within an environment, or find the locations of particular molecules and structures within a cell. Approaches have also been developed to distinguish living from dead cells based on whether they take up particular fluorochromes. Sometimes, multiple fluorochromes are used on the same specimen to show different structures or features.

One of the most important applications of fluorescence microscopy is immunofluorescence, which is used to identify certain microbes by observing whether antibodies bind to them. (Antibodies are protein molecules the immune system produces that attach to specific pathogens to kill or inhibit them.) This technique has two approaches: direct immunofluorescence assay (DFA) and indirect immunofluorescence assay (IFA). In DFA, specific antibodies (e.g., those that target the rabies virus) are stained with a fluorochrome. If the specimen contains the targeted pathogen, one can observe the antibodies binding to the pathogen under the fluorescent microscope. This is a primary antibody stain because the stained antibodies attach directly to the pathogen.

In IFA, secondary antibodies are stained with a fluorochrome rather than primary antibodies. Secondary antibodies do not attach directly to the organism but bind to primary antibodies. When the unstained primary antibodies bind to the pathogen, the fluorescent secondary antibodies can be observed binding to the primary antibodies. Thus, the secondary antibodies are attached indirectly to the pathogen. Since multiple secondary antibodies can often attach to a primary antibody, IFA increases the number of fluorescent antibodies attached to the specimen, making it easier to visualize its features.

This text is adapted from Openstax, Microbiology 2e, Section 2.4: Staining Microscopic Specimens.

 Core: Cell Biology

Unrenewable Cells

JoVE 13472

In humans, the photoreceptor cells of the eye and sensory hair cells of the ear lack stem cells. These cells are thus unrenewable and cannot be replaced when they are damaged or destroyed.

Photoreceptors

The retina is composed of several layers and contains specialized cells called photoreceptors. The photoreceptors (rods and cones) change their membrane potential when stimulated by light energy. There are two types of photoreceptors—rods and cones—which differ in the shape of their outer segment. The rod-shaped outer segments of the rod photoreceptor contain a stack of membrane-bound discs that contain the photosensitive pigment rhodopsin. The cone-shaped outer segments of the cone photoreceptor contain their photosensitive pigments in infoldings of the cell membrane. Since the eye photoreceptors lack stem cells, the loss of rods and cones is permanent and cannot be regenerated. The loss of photoreceptors or damage can be caused by age-related problems or exposure to high-intensity light.

Sensory Hair Cells

Like photoreceptors, the ear's sensory hair cells also do not contain stem cells. The hair cells are present in the organs of Corti, which are named for the hair-like stereocilia extending from the cell’s apical surfaces. When the pressure waves from the scala move the basilar membrane, the tectorial membrane slides across the stereocilia. This bends the stereocilia either toward or away from the tallest member of each array. When the stereocilia bend toward the tallest member of their array, the tension in the protein tethers opens ion channels in the hair cell membrane. This will depolarize the hair cell membrane, triggering nerve impulses that travel down the afferent nerve fibers attached to the hair cells. When the stereocilia bend toward the shortest member of their array, the tension on the tethers slackens, and the ion channels close. When no sound is present and the stereocilia are standing straight, a small amount of tension still exists on the tethers, keeping the membrane potential of the hair cell slightly depolarized. The loss or damage of hair cells is also permanent and often leads to permanent deafness.

This text is adapted from Openstax Anatomy and Physiology 2e, Section 14.1: Sensory Perception.

 Core: Cell Biology

Handwashing III: During the Procedure and Post-Procedure Steps

JoVE 13653

To wash hands properly, follow these steps:

  1. Wet the hands. Use enough soap to cover all surfaces of the hands.
  2. Rub both hands palm to palm.
  3. Rub the back of the hand: Use the right palm over the left dorsum or back of the left hand with interlaced fingers, then switch hands.
  4. Rub palm to palm with fingers interlaced.
  5. With fingers interlocked, rub the backs of the fingers of the opposite hands. 
  6. Rub the left thumb clasped in the right palm in a circular motion and repeat with the other thumb.
  7. Rub the tips of the fingers. Rotationally rub backward and forward with clasped fingers of the right hand in the left palm and vice versa.
  8. Rub each wrist with the opposite hand.
  9. Rinse hands with water, as soap residue can lead to irritation and damage to the skin. Damaged skin does not provide a barrier to infection for the healthcare worker and can become colonized with potentially pathogenic bacteria, leading to cross-infection in the patient.
  10. Ensure all surfaces of the hands are cleaned. Areas that are missed can be a source of cross-infection as well.
  11. Following hand washing, turn off the faucet with a paper towel or use the elbow to prevent contact with the faucet. Dispose of paper towels.
  12. With another towel, dry the hands. Wet hands encourage the growth of bacteria. Dispose of used paper towels in a foot-operated waste bin to prevent contamination of the hands.

 Core: Nursing

Types of Errors: Detection and Minimization

JoVE 14506

Error is the deviation of the obtained result from the true, expected value or the estimated central value. Errors are expressed in absolute or relative terms.

Absolute error in a measurement is the numerical difference from the true or central value. Relative error is the ratio between absolute error and the true or central value, expressed as a percentage.

Errors can be classified by source, magnitude, and sign. There are three types of errors: systematic, random, and gross.

Systematic or determinate errors emerge from known sources and are reproducible during replicate measurements. Defective equipment and experiment design flaws are familiar sources of these errors. These errors can be minimized by employing standard reference materials, independent analysis, or varying the sample size.

Random errors or indeterminate errors are difficult to reproduce by repeating measurements. These errors originate from uncontrolled variables like electronic noise in the circuit of an electrical instrument and irregular changes in the heat loss rate from a solar collector due to changes in the wind.

Gross errors are caused by human mistakes. The magnitude of these errors is often high. The origin of such errors is entirely based on the observer.

 Core: Analytical Chemistry

Quantifying and Rejecting Outliers: The Grubbs Test

JoVE 14522

Sometimes, a data set can have a recorded numerical observation that greatly  deviates from the rest of the data. Assuming that the data is normally distributed, a statistical method called the Grubbs test can be used to determine whether the observation is truly an outlier.  To perform a two-tailed Grubbs test, first, calculate the absolute difference between the outlier and the mean. Then, calculate the ratio between this difference and the standard deviation of the sample. This number is known as the Grubbs statistic, 'G.' When the calculated G value exceeds the G critical value for a given confidence level and the number of observations, the questionable observation is considered an outlier and removed from the data set. On the contrary, if the calculated G value is smaller than the critical G value, the questionable observation is not considered an outlier and therefore retained in the data set.

 Core: Analytical Chemistry

Redox Equilibria: Overview

JoVE 14538

A reduction-oxidation reaction is commonly called a redox reaction. In a redox reaction, electrons are transferred from one species to another rather than being shared between or among atoms. The reducing agent or reductant is the species that loses electrons and gets oxidized in the process. The species that gains electrons and gets reduced in the process is the oxidizing agent or oxidant. Redox reactions are represented as two separate equations called half-reactions, where one equation represents the species that gains electrons, and the other represents the species that loses electrons.

Generally, the thermodynamics of a reaction is expressed in terms of the change in Gibbs free energy (ΔG), which is a function of concentrations of reactants and products. However, the thermodynamics of a redox reaction is expressed in terms of electrochemical potential (E)  and the Nernst equation, as the reaction involves the movement of electrons. The Nernst Equation expresses the relationship between the electrochemical potential and the concentrations of the reactants and products.

 Core: Analytical Chemistry

Effects of EDTA on End-Point Detection Methods

JoVE 14577

Different methods, such as visual observance of metal-ion indicators, spectroscopic techniques, and potentiometric methods, can determine the endpoint of an EDTA titration.

In the visual method, metal-ion indicators (metallochromic dyes), which have distinct colors in their free and complex forms, are added to the mixture to signal the titration's end point. They form stable complexes with metal ions, but these complexes are weaker than the corresponding metal–EDTA complexes. As a result, EDTA will displace the indicator to form a more stable metal–EDTA complex as the endpoint is reached.

Since the color of the free indicators also depends on pH, titration must be performed in the appropriate pH range.

If the titration solution is already intensely colored or the endpoint color change is subtle, the visual method might not be precise enough. In such cases, the spectroscopic method is employed. Here, the endpoint is detected by monitoring the absorbance of the mixture at a particular wavelength.

Another method is potentiometric titration, which detects the endpoint by measuring the change in the potential of the metal ion during the titration using an appropriate electrode for the metal.

 Core: Analytical Chemistry

Tandem Mass Spectrometry

JoVE 14593

Tandem mass spectrometry is a technique that uses multiple mass analyzers in series to obtain a higher selectivity and signal-to-noise ratio for the analyte. Instruments with multiple analyzers separated by an interaction cell enable secondary fragmentation and selected study of the fragment ions.

Secondary fragmentations occur in the interaction cell and can be induced by various factors. Fragmentation induced by collision with inert gases, such as N2, Ar, He, etc., is called collision-induced dissociation (CID). Other techniques involve interaction with the surface through surface-induced dissociation (SID); interaction with low-energy electrons, called electron-capture dissociation (ECD); or through an intense laser beam, called photo-induced dissociation (PID).

Tandem mass spectrometry can be used for various experiments by varying the mode of operation of the two mass analyzers. When the first mass analyzer selects a particular precursor ion and the second mass analyzer scans all the fragmented product ions, it is called a product-ion scan, and the spectra obtained by this experiment are called product-ion spectra. Alternately, when the first mass analyzer scans all the precursor ions and the second mass analyzer selects a particular product ion, it is called 'precursor-ion mode', and the spectra are termed precursor-ion spectra. Another experiment can be performed by scanning both the analyzers simultaneously with an offset of a mass equal to a neutral fragment lost by the ion to be monitored. Such a mode is called 'neutral-loss scanning mode', and the corresponding spectra are called neutral-loss spectra.

The instrumentation of tandem mass spectrometry can be varied. Tandem-in-space spectrometers use two independent mass analyzers in two different regions in space, for example, a triple quadrupole mass spectrometer and TOF–TOF spectrometer. On the other hand, a tandem-in-time spectrometer, such as a quadrupole ion-trap instrument, uses the same spatial region to fragment the ions and dissociate and analyze a selected ion of interest, but at different times by expelling the unwanted ions. Tandem-in-time spectrometers can be used to perform MSn experiments.   

 Core: Analytical Chemistry

EDTA: Indirect and Alkalimetric Titration

JoVE 14730

Unlike direct titration, back-titration, and displacement titration, indirect titration is an EDTA titration method for quantifying anions. In the indirect titration method, anions are precipitated as their insoluble salts with excess metal ions. The filtrate containing the excess metal ions is directly titrated with standard EDTA until the endpoint is achieved. Another approach involves extracting the metal ion and back-titrating with standard EDTA to obtain the endpoint. In this way, the amount of anion is indirectly estimated.

Another EDTA titration method is alkalimetric titration. This titration method is similar to acid–base titration. During the titration, disodium EDTA is added to a metal ion solution, forming the metal–EDTA complex. This process involves the liberation of two equivalents of hydrogen ions in the solution, which is titrated using a base such as standard sodium hydroxide in the presence of appropriate acid–base indicator until the end-point concentration is obtained.

 Core: Analytical Chemistry

Structure and Organization of Smooth Muscles

JoVE 14855

Smooth muscle tissue is a type of muscle tissue that can be found lining various vital organs in the human body, including the lungs, blood vessels, digestive tract, and respiratory tract. This type of tissue is responsible for regulating the movements of these organs, playing crucial roles in the functioning of various systems, including the vascular, digestive, respiratory, and urinary systems.

Structure of smooth muscle cell

Smooth muscle cells are spindle-shaped with tapering ends and a single nucleus in the center. They typically range from 30 to 200 µm in length and have a diameter of about 5-10 µm. Unlike the striated appearance of skeletal and cardiac muscles, smooth muscle fibers have a uniform, non-striated appearance. This appearance is due to the lack of organized sarcomeres. The contractile proteins actin and myosin, which form the thin and thick filaments, respectively, are arranged obliquely and irregularly along the long axis of the cell. The thin filaments are anchored to the sarcolemma by dense bodies interconnected via the intermediate filament network. Instead of T-tubules, these muscle fibers have caveolae — small pockets of invaginated sarcolemma that help influx Ca2+ ions during contraction cycles.

Smooth muscles can be further categorized into two types: multi-unit and single-unit (visceral) smooth muscle.

Multi-unit smooth muscle

In contrast, multi-unit smooth muscle fibers consist of individual muscle fibers that operate more independently of each other. Each fiber has its own nerve supply, and there are fewer gap junctions between cells. This arrangement allows for more precise and localized control of contraction. Multiunit smooth muscle is found in areas where fine control is necessary, such as in the iris of the eye, where it controls pupil size, and in the walls of blood vessels, where it regulates blood flow and pressure.

Visceral smooth muscle

Visceral smooth muscle is the most common type of smooth muscle. It is found in the walls of hollow organs such as the stomach, intestines, uterus, and bladder. These cells are connected by gap junctions, which allow electrical and chemical signals to pass directly from one cell to another. As a result, when one cell becomes excited and contracts, the signal quickly spreads to neighboring cells, causing the muscle layer to contract as a synchronized unit. This property is essential for functions like peristalsis in the gastrointestinal tract, where coordinated waves of contraction move contents through the system.

 Core: Anatomy and Physiology

Muscles that Move the Arm

JoVE 14876

Nine muscles are involved in arm movements. Two of these, the pectoralis major and latissimus dorsi, originate from the axial skeleton and are called axial muscles. The other seven originate from the scapula and are called the scapular muscles.

The pectoralis major has two origins. Its clavicular head originates on the medial half of the clavicle. In contrast, the sternocostal head originates on the costal cartilages of ribs 1-6, the sternum, and the aponeurosis of the external oblique of the anterior abdominal wall. Both these heads merge and insert on the intertubercular sulcus of the humerus. Its clavicular head flexes the humerus, and the sternocostal head extends a flexed arm back to the anatomical position. When the head works together, they adduct and medially rotate the humerus.

The latissimus dorsi originates from the spinous processes of the thoracic vertebrae T7 to T12, crests of sacrum and ilium, and ribs 9 to 12. It inserts anteriorly on the intertubercular sulcus of the humerus. Acting with the scapular muscle, teres major, it extends, adducts, and medially rotates the humerus. The origin of the scapular muscle, teres major, is at the inferior angle of the scapula. It inserts into the humerus at the medial lip of its intertubercular sulcus. Its actions are the same as that of latissimus dorsi.

The scapular deltoid is the most superficial muscle of the shoulder joint and is responsible for its rounded shape. Its anterior fibers originate on the lateral of the clavicle — the lateral fibers originate from the acromion of the scapula, and the posterior fibers originate from the spine of the scapula. All the heads merge and insert on the deltoid tuberosity of the humerus. While the lateral fibers abduct the arm, the anterior fibers flex and medially rotate it, and the posterior fibers extend and laterally rotate it.

The scapular muscles — supraspinatus, infraspinatus, teres minor, subscapularis, and their corresponding tendons, combine to create the rotator cuff. The origin of the supraspinatus muscle lies on the supraspinous fossa of the scapula. Conversely, the infraspinatus muscle originates from its infraspinous fossa. Both these muscles are inserted at the greater tubercle of the humerus, aiding in arm abduction and lateral rotation, respectively.

Another muscle, the teres minor, also inserts at the greater tubercle of the humerus but originates from the inferior angle of the scapula. It acts along with the infraspinatus muscle to laterally rotate and extend the arm. The origin of the subscapularis muscle is at the subscapular fossa of the scapula. It inserts on the lesser tubercle of the humerus and aids in the medial rotation of the arm. Finally, the last scapular muscle involved in arm movements is the coracobrachialis. The origin of this muscle is at the coracoid process of the scapula, and the insertion is on the middle surface of the medial shaft of the humerus. It helps with arm flexion and adduction.

The rotator cuff muscles are prone to injuries, especially in people active in sports. Due to its positioning between the head of the humerus and the acromion of the scapula, the supraspinatus muscle is particularly prone to strain and degenerative changes.

 Core: Anatomy and Physiology

Action Potential: Phases of Stimulation

JoVE 14893

The action potential is a complex electrical event that occurs in excitable cells, such as neurons and muscle cells. It consists of several distinct phases, each with specific characteristics.

Resting Phase:

In this phase, the cell's membrane is at its resting potential, typically around -70 millivolts (mV) for neurons. Inside the cell, there is a higher concentration of potassium ions (K+) and a lower concentration of sodium ions (Na+). Voltage-gated sodium channels are closed, and voltage-gated potassium channels are closed but capable of opening.

Depolarization Phase:

A graded potential, often an excitatory postsynaptic potential (EPSP), reaches the threshold level (typically around -55 mV). This triggers the voltage-gated sodium channels to open rapidly, allowing an influx of sodium ions into the cell. This rapid sodium influx causes a sharp increase in membrane potential, turning it more positive. The influx of sodium ions further depolarizes the membrane, leading to a positive feedback loop that triggers more sodium channels to open.

Peak of the Action Potential:

At the peak of the action potential, the sodium channels begin to inactivate or close, reducing sodium influx. Voltage-gated potassium channels start to open slowly in response to the increasing membrane potential.

Repolarization Phase:

As voltage-gated potassium channels open fully, potassium ions exit the cell. This movement of positively charged ions out of the cell helps to restore the negative membrane potential. The membrane potential gradually returns to the resting potential of around -70 mV.

Hyperpolarization Phase (Undershoot):

The movement of potassium ions continues for a brief period, causing the membrane potential to dip below the resting potential, typically around -80 mV. The delayed closure of some potassium channels contributes to this temporary hyperpolarization.

Refractory Period:

During and immediately after an action potential, it is impossible to trigger another one. This prevents the action potential from moving backward. This is called the absolute refractory period.

Following the absolute refractory period, it is possible to initiate another action potential, but it requires a stronger stimulus than usual. This is known as the relative refractory period.

The phases of an action potential are essential for transmitting electrical signals in neurons. This rapid and coordinated sequence of events allows for the unidirectional propagation of signals along the length of the neuron, enabling communication within the nervous system and with other cells.

 Core: Anatomy and Physiology

Diencephalon: Anatomical Regions

JoVE 14909

The diencephalon, etymologically translated as 'through brain,' plays an integral role as the conduit between the cerebrum and the vast extent of the nervous system. However, the olfactory system is an exception, as it interfaces directly with the cerebrum. The diencephalon, deeply ensconced beneath the cerebrum, primarily consists of three paired structures — the thalamus, hypothalamus, and epithelamus. It also includes accessory structures such as the subthalamus, which houses the subthalamic nucleus as part of the basal nuclei.

Thalamus

The thalamus operates as a central processing unit, encompassing various nuclei that transmit information between the cerebral cortex, spinal cord, and brainstem. All sensory information, except olfaction, is routed through the thalamus for preprocessing before it reaches the cortex. Axons from the peripheral sensory organs, or intermediary nuclei, synapse in the thalamus, and thalamocortical neurons project directly to the cerebrum. It is essential to note that the thalamus is more than a passive courier; it actively processes incoming information. For instance, the segment of the thalamus receiving visual data will dictate which stimuli warrant attention.

The thalamus also receives information from the cerebellum, typically in the form of motor commands. This interaction involves crosstalk with the cerebellum and other brain stem nuclei. An intricate network of connections exists between the cortex and basal nuclei involving the thalamus. The neuronal information relayed by the basal nuclei is directed toward the thalamus, which subsequently relays this output to the cerebral cortex. The cortex, in turn, transmits data to the thalamus, which influences the activity of the basal nuclei.

Hypothalamus

The hypothalamus, another primary region of the diencephalon, is positioned anterior and slightly below the thalamus. The hypothalamus includes a series of nuclei primarily tasked with maintaining homeostasis. This region serves as the command center governing the autonomic nervous system and the endocrine system via regulation of the anterior pituitary gland. Additional sections of the hypothalamus participate in memory and emotional processes as constituents of the limbic system.

Epithalamus

The epithalamus, the smallest part of the diencephalon, is situated posterior to the thalamus. This area contains the habenular nuclei, which relay information from the limbic system to the midbrain. The habenula is also thought to influence reward and punishment processing behavior. Furthermore, the epithalamus houses the pineal gland, which helps regulate sleep-wake cycles through melatonin release. It is one of the few parts of the brain thought to maintain some degree of neuroplasticity throughout life. The epithalamus is an essential part of the circuitry connecting the limbic system to motor regions, and it plays a key role in integrating sensory input from other areas. Additionally, it helps control endocrine functions by regulating hormones produced by the pituitary gland. Research also suggests that the epithalamus plays a role in emotion, learning, and memory formation.

 Core: Anatomy and Physiology

Spinal Nerves: Plexus I

JoVE 14925

Nerve plexuses are networks of interlacing nerves that serve as communication hubs to distribute and organize nerve action across various body regions. The nerve plexuses are organized into the cervical plexus located in the neck region, brachial plexus in the shoulder area, lumbar plexus found in the lower back, sacral plexus situated in the pelvis, and coccygeal plexus located in the coccygeal region.

The Cervical Plexus

The cervical plexus, formed by the anterior rami of the first four cervical spinal nerves (C1-C4), and partially the fifth cervical spinal nerve (C5), is situated in the neck region. This plexus innervates the skin and muscles of the head, neck, and shoulders, facilitating both sensory and motor functions.

  • • Motor Innervation: The phrenic nerve is a primary motor nerve emerging from the cervical plexus. It innervates the diaphragm and is essential in controlling the respiratory process. Without the proper function of this nerve, breathing would be significantly impaired.

Other motor components of the cervical plexus include the ansa cervicalis and segmental branches. The ansa cervicalis innervates some infrahyoid muscles, aiding in swallowing and speech. The segmental branches supply motor fibers to the deep muscles of the neck, contributing to the stability and movement of the head.

  • • Sensory Innervation: The cervical plexus also comprises four primary sensory nerves: the lesser occipital, great auricular, transverse cervical, and supraclavicular nerves. These nerves provide sensation to the skin of the neck, ear area, and parts of the shoulder and chest, ensuring sensory input from these regions reaches the central nervous system.

The Brachial Plexus

The brachial plexus extends from the neck into the axilla and is formed by the anterior rami of the fifth to eighth cervical nerves (C5-C8) and the first thoracic nerve (T1). This plexus supplies the shoulder and upper limbs, orchestrating a wide range of movements and sensory functions. The brachial plexus gives rise to five significant nerves:

  • • Axillary Nerve: This innervates the deltoid and teres minor muscles, enabling shoulder abduction and playing a role in the sensation of the shoulder area.
  • • Musculocutaneous Nerve: This nerve supplies the flexor muscles in the front of the arm, facilitating elbow flexion, and provides sensory information from the lateral forearm.
  • • Radial Nerve: By innervating the muscles on the posterior aspect of the upper limb, this nerve supports elbow extension, wrist and finger extension, and thumb abduction. It also conveys sensory input from the posterior arm and hand.
  • • Median Nerve: This nerve is primarily responsible for flexing the wrist and fingers, along with thumb opposition. It also carries sensory information from the palmar aspect of the hand.
  • • Ulnar Nerve: The intricate movements of the fingers and wrist are controlled by this nerve. It provides sensation to the little finger and part of the ring finger, highlighting its importance in grip and hand movements.

 Core: Anatomy and Physiology

Higher Mental Functions of the Brain: Learning and Memory

JoVE 14944

Memory is one of the most vital higher mental functions of the brain. Memory is closely related to learning because it enables us to retain information and experiences from our past to use them in our present life. It also helps us to remember facts, events, and skills, such as riding a bike or swimming. There are two types of memory — declarative memory, which involves memorizing facts or events, and procedural memory, which enables us to remember how to do something like writing or playing an instrument.

Memory can be further divided into short-term and long-term memory. Short-term memory temporarily stores information, while long-term memory stores the information for longer. The hypothalamus consolidates short-term memories in our brain, which means that the nervous system will transfer information from the hypothalamus to the cerebral cortex for permanent storage. This means that consolidated short-term memories get converted into long-term memories.

Amnesia, characterized by forgetfulness, refers to the absence or impairment of memory. It can manifest as a complete or partial inability to recall past experiences. Anterograde amnesia involves explicitly the loss of memory for events that transpire after the underlying trauma or disease, failing to form new memories. On the other hand, retrograde amnesia pertains to the loss of memory for events before the trauma or disease, leading to an inability to recollect past events.

 Core: Anatomy and Physiology

Focusing of Light in the Eye

JoVE 14963

Light rays enter the eye through the cornea, a transparent dome-shaped tissue that is the eye's outermost layer. The cornea bends or refracts, light rays traveling to the pupil. The shape of the cornea determines how much of the light is bent and whether the image will be focused correctly on the retina at the back of the eye. Once the light has passed through both refraction layers, it converges into a single focal point onto a small area. This is where photoreceptors start transforming visible photons into electrical signals sent along nerve fibers up to your brain for interpretation - ultimately resulting in what is known as sight.

The refractive power of the human eye is its ability to bend light rays as they enter the eye to focus them onto the retina at the back of the eye. This is necessary because light entering the eye is initially divergent, but the images we perceive must be clear and focused. The refractive power of the human eye is measured in diopters (D).

The cornea is the primary refractive surface of the eye, providing approximately two-thirds of the eye's refractive power. The lens of the eye, located behind the iris, provides the remaining one-third of the eye's refractive power.

The lens can change shape, known as accommodation, allowing the eye to focus on objects at different distances. The eye's crystalline lens is located directly behind the iris and pupil and further focuses light onto the retina. Unlike a rigid camera lens, our crystalline lenses are elastic and can change shape depending on one's vision needs for near or far objects.

When looking at an object closer than 20 feet away from you (near vision), your eyes must accommodate by becoming more curved to focus on near objects properly. When looking at an object farther than 20 feet away from you (distant vision), your eyes must become less curved to focus on distant objects properly.

The amount of bend or curve placed upon incoming light rays depends on one's refractive power index: The higher power index means more curvature required of incoming light, and vice versa. People with myopia (nearsightedness) have too much anatomical curvature in their eyes. This causes incoming light to focus too soon before reaching the retina, causing blurry vision when looking at distant objects. Inversely, people with hyperopia (farsightedness) have anatomical curvature that is too slight in their eyes. This causes incoming light rays to not bend enough before reaching their retinas and, as a result, causes blurry vision when looking at nearby objects.

Some other types of refractive errors can occur in the human eye, including:

  1. Astigmatism occurs when the cornea, the eye's clear front surface, is shaped like a football instead of spherical, causing blurry and distorted vision at any distance.
  2. Presbyopia is an age-related condition that typically occurs when the eye lens is less flexible, resulting in difficulty focusing on close-up objects.

These refractive errors can be corrected using appropriate glasses or contact lenses or surgery in some extreme cases.

 Core: Anatomy and Physiology

The Pituitary Gland

JoVE 14979

The pituitary is a small endocrine organ in the sphenoid bone under the hypothalamus. Primarily, the pituitary in adults has two distinct anatomical and functional regions— the anterior and posterior lobes. During human fetal development, a third pituitary gland region called the pars intermedia atrophies and disappears. However, some of its cells migrate and exist adjacent to the anterior pituitary in adults.

  • The anterior lobe comprises the pars distalis and the pars tuberalis. These areas consist of glandular epithelial tissue that produces various hormones. The anterior lobe is linked to the hypothalamus through the hypophyseal portal system. The hormones secreted by the ventral hypothalamus cells travel through this portal system and regulate hormone-secreting cells of the anterior lobe.

The posterior pituitary lobe has two parts: the infundibulum, a funnel-shaped stalk, and the pars nervosa. The infundibulum allows the passage of the supra-optic and paraventricular nuclei axons of the hypothalamic-hypophyseal tract, maintaining a neural connection with the hypothalamus. The hypothalamic-hypophyseal tract delivers hormones secreted by hypothalamic neurosecretory cells into the posterior pituitary lobe for storage and release.

 Core: Anatomy and Physiology

Voltage

JoVE 15064

The movement of electrons in a conductor requires some form of energy or work, usually provided by an external force, like a battery. This force is called the electromotive force or voltage. The voltage between two points, referred to as points "a" and "b," in an electric circuit is the energy (or work) needed to move a unit charge from point "a" to point "b," and this relationship is expressed mathematically as

Equation1

In this equation, "w" represents the energy measured in joules (J), and "q" represents the charge measured in coulombs (C). The voltage, denoted as "vab" is measured in volts (V).

Voltage, often referred to as the potential difference, signifies the energy required to move a unit charge through an element within the circuit. It is important to note that the value of a voltage can be either positive or negative, with its direction determined by its polarities, indicated as (+) and (-). In electrical terminology, it is customary to state that a voltage exists across an element. The notation "vab" represents the voltage between points "a" and "b" and can be interpreted in two distinct ways:

  • • Point "a" is at a potential of "vab" volts higher than point "b."
  • • The potential at point "a" with respect to point "b" is "vab" volts.

A constant voltage is categorized as a direct current (DC) voltage, typically represented as "V." DC voltages are commonly generated by sources such as batteries. On the other hand, a voltage that varies sinusoidally with time is called an alternating current (AC) voltage, represented as "v." AC voltages are typically produced by electric generators.

 Core: Electrical Engineering

Characteristics of Practical Op Amps

JoVE 15080

A difference amplifier, a crucial component in numerous electronic devices, ideally amplifies only the difference-mode signal, which is the difference between two input signals. However, in practical circuits, the output voltage depends on both the differential gain and the common-mode gain.

The ratio of differential gain to the common-mode gain is defined as the common-mode rejection ratio (CMRR). This ratio quantifies the ability of operational amplifiers (op-amps) to reject common-mode signals, which are identical signals present at both inputs of the amplifier.

For instance, consider a sound system where an op-amp picks up both music and electrical noise. A high CMRR allows the op-amp to effectively ignore the noise and amplify only the music, enhancing the sound quality of the system.

Another crucial characteristic of op-amps is the gain-bandwidth product. This is the product of the op-amp's open-loop voltage gain and the frequency at which this gain is measured. It remains constant for any given op-amp, providing a useful parameter for comparing the performance of different op-amps.

A unity gain buffer, or voltage follower, is a specific configuration of an op-amp that provides a voltage gain of one. This means that the output voltage matches the input voltage. It has low output impedance and high input impedance, making it ideal for signal isolation and impedance matching.

In an audio system, a voltage follower plays a vital role in preventing a low-impedance speaker from effectively shorting the audio signal from a high-impedance source to the ground. By matching the impedance levels, it ensures that maximum power is transferred from the source to the speaker without distortion or loss.

In conclusion, understanding the characteristics and functionalities of op-amps, such as the CMRR, gain-bandwidth product, and unity gain buffer configuration, is essential in designing effective and efficient electronic systems. These principles form the backbone of many modern audio and communication systems, influencing their performance and quality.

 Core: Electrical Engineering

Second-Order Circuits

JoVE 15097

Integrating two fundamental energy storage elements in electrical circuits results in second-order circuits, encompassing RLC circuits and circuits with dual capacitors or inductors (RC and RL circuits). Second-order circuits are identified by second-order differential equations that link input and output signals.

Input signals typically originate from voltage or current sources, with the output often representing voltage across the capacitor and/or current through the inductor. For example, in an RLC circuit, initial energy stored in the capacitor and inductor initiates the circuit. Applying Kirchhoff's voltage law and performing a time derivative yields a second-order differential equation.

Equation1

 Its coefficients, determined by resistance, capacitance, and inductance, manifest as the damping coefficient and resonant frequency.

Equation2

Equation3

The damping coefficient plays a critical role in these circuits, signifying the extent of damping caused primarily by the resistor. It directly influences the pace at which energy dissipates within the system, effectively controlling the rate of energy loss. On the other hand, the resonant frequency is a key characteristic that represents the circuit's innate oscillation frequency. It measures how quickly energy is exchanged between the inductor and capacitor in the circuit, illustrating the circuit's natural tendency to oscillate at a particular frequency.

The damping coefficient dictates how fast energy is lost in the system due to resistance. At the same time, the resonant frequency highlights the circuit's inherent oscillation speed as energy shifts between the inductor and capacitor. These two factors are crucial in understanding and analyzing the behavior of second-order circuits.

 Core: Electrical Engineering

Norton Equivalent Circuits

JoVE 15114

Norton's theorem is a fundamental concept in the field of electrical engineering that allows for the simplification of complex AC circuits. The theorem states that any two-terminal linear network can be replaced with an equivalent circuit that consists of an impedance, which is parallel with a constant current source. Figure 1 shows the AC circuit portioned into two parts: Circuit A and Circuit B, while Figure 2 depicts the circuit obtained by replacing Circuit A by its Norton equivalent circuit.

Figure1

Figure 1: A circuit portioned into two parts 

Figure2

Figure 2: Norton equivalent circuit

To calculate the value of the parallel impedance, one must replace the source with its internal impedance, resulting in a circuit with an equivalent impedance known as the Norton impedance. The Norton impedance is the same as the Thévenin impedance and is used to determine the Norton current, which is the current flowing through the circuit.

Determining the Norton current requires placing the sources back into the circuit and analyzing the open-circuit voltage, also known as the Thévenin voltage. The value of the Thévenin voltage is determined by multiplying the source current by the Thevenin impedance and is used to drop the same voltage across the load impedance when it is placed in a parallel configuration.

By using the relationship between the Norton current, the Thévenin voltage, and the Norton current values, one can determine the Norton current of the circuit. This relationship makes Norton's theorem beneficial for analyzing and designing systems containing complex AC circuits since it simplifies their analysis by breaking the circuit down into smaller, more manageable sections.

 Core: Electrical Engineering

Mixtures of Acids

JoVE 17362

The pH of a solution containing an acid can be determined using its acid dissociation constant and initial concentration. If a solution contains two different acids, then its pH can be determined using one of several methods depending on the relative strength of the acids and their dissociation constants.

In a strong and weak acid mixture, the strong acid dissociates completely and becomes a source of almost all the hydronium ions present in the solution. In contrast, the weak acid shows partial dissociation and produces a negligible concentration of hydronium ions. The high concentration of hydronium ions produced by the strong acid further reduces the dissociation of the weak acid. According to Le Chatelier's principle, this happens: "When a chemical system at equilibrium is disturbed, the system shifts in a direction that minimizes the disturbance." The excess hydronium ions produced by the strong acid disturb the equilibrium, and thus, the reaction will move in the reverse direction until the equilibrium is established. This leads to a decrease in the dissociation of the weak acid. Because of this decrease, the pH of a strong and weak acid mixture can be calculated from the concentration of the strong acid only. For example, the pH of a mixture with an equal concentration of hydrochloric acid (HCl), a strong acid, and formic acid (HCHO2), a weak acid, can be determined from the concentration of HCl only. If the concentration of the HCl in the mixture is 0.0020 M, its pH can be calculated as follows.

pH = −log(0.002) = 2.7

Here, the concentration of hydronium ions produced by HCHO2 and the autoionization of water is negligible and thus can be ignored.

A Mixture of Two Weak Acids with Different Dissociation Constants

In a mixture of two weak acids, the pH of a mixture will be determined by the stronger acid if its dissociation constant is significantly higher than the weaker acid. For example, in a mixture with an equal concentration of nitrous acid (HNO2) and hypochlorous acid (HClO), the HNO2 will be the main determinant of the pH of the mixture as its Ka (4.6 × 10−4) is approximately 10,000 times higher than the Ka (2.9 × 10−8) of HClO. According to Le Chatelier's principle, HClO shows decreased dissociation in the presence of HNO2.

 Core: Analytical Chemistry

Isolating Nucleic Acids from Yeast

JoVE 5096

One of the many advantages to using yeast as a model system is that large quantities of biomacromolecules, including nucleic acids (DNA and RNA), can be purified from the cultured cells.

This video will address the steps required to carry out nucleic acid extraction. We will begin by briefly outlining the growth and harvest, and lysis of yeast cells, which are the initial steps common to the isolation of all biomacromolecules. Next, we will discuss two unique purification methods for the separation of nucleic acids: column binding and phase separation. Additionally, we will demonstrate several ways in which these methods are applied in the laboratory, including the preparation of nucleic acids for molecular biology techniques such as PCR and southern blotting, quantification of gene expression in response to environmental stimuli, and purification of large amounts of recombinant proteins.

 Biology I

Basic Mouse Care and Maintenance

JoVE 5158

Mice (Mus musculus) are small rodents that breed and sexually mature quickly, making them perfectly suited to generating large animal colonies for biological research. As compared to other mammalian species, mice are simple and inexpensive to maintain in the laboratory. Nevertheless, mouse colonies do have specific husbandry needs that are critical to preserving animal health and safety as well as experimental reproducibility.

This video demonstrates standard practices that ensure mice are treated as humanely as possible within the laboratory animal facility, or vivarium. The discussion begins by reviewing a typical mouse housing setup, consisting of a plastic cage equipped with a layer of soft bedding and nesting material. The preformulated food pellets (also known as chow) that comprise the typical mouse diet are also introduced. In order to facilitate experiments performed on mice, safe animal handling practices are demonstrated, including common restraint techniques like “scruffing,” and the strategies used by researchers to keep track of individual mice within the facility. Finally, experimental manipulations of mouse housing and diet are discussed, in addition to one of the most common applications of the scruffing technique — performing injections.

 Biology II

The Morris Water Maze

JoVE 5211

The Morris water maze is one of the most widely used behavioral tests for studying spatial learning and memory. In the initial phases of this task, rodents must swim to a platform to escape from a pool of water. The platform is then hidden under the water’s surface, so that the animal is required to remember it’s location in order to escape.This simple yet powerful maze design can be used to assay cognitive function, study animal models of neurodegenerative disease, and test potential drug therapies.

This video provides an introduction to the Morris water maze and the principles surrounding its use, including a discussion of the different types of memory tested in the maze, important points to consider when designing and conducting this experiment, and the procedures for setup and running of the test. Several applications of the maze are examined, such as investigating how radiation treatment may lead to memory impairment. Finally, other types of water mazes, such as the 8-arm radial maze, are introduced to show how this paradigm can be adapted to engage different types of memory.

 Neuroscience

Induced Pluripotency

JoVE 5333

Induced pluripotent stem cells (iPSCs) are somatic cells that have been genetically reprogrammed to form undifferentiated stem cells. Like embryonic stem cells, iPSCs can be grown in culture conditions that promote differentiation into different cell types. Thus, iPSCs may provide a potentially unlimited source of any human cell type, which is a major breakthrough in the field of regenerative medicine. However, more research into the derivation and differentiation of iPSCs is still needed to actually use these cells in clinical practice.

This video first introduces the fundamental principles behind cellular reprogramming, and then demonstrates a protocol for the generation of iPSCs from differentiated mouse embryonic fibroblasts. Finally, it will discuss several experiments in which scientists are improving or applying iPSC generation techniques.

 Developmental Biology

Balance and Coordination Testing

JoVE 5423

Balance and coordination are critical components involved in the control of movement. Many sensory receptors and neural processing units are required to help individuals maintain balance while performing various activities. Deficits in balance and coordination occur in patients suffering from movement disorders or due to aging. Therefore, scientists are trying to understand the pathophysiology behind these conditions. One way to do that is by using rodent models and testing them on behavioral paradigms such as the rotarod or balance beam.

This video discusses the currently known neurophysiology behind balance and coordination. Then, we go over protocols to run balance tests in rodents using the rotarod and balance beam. Finally, we'll discuss some current studies utilizing these methods to investigate aging, muscular dystrophy and Parkinson's disease.

 Behavioral Science

Solid-Liquid Extraction

JoVE 5538

Source: Laboratory of Dr. Jay Deiner — City University of New York

Extraction is a crucial step in most chemical analyses. It entails removing the analyte from its sample matrix and passing it into the phase required for spectroscopic or chromatographic identification and quantification. When the sample is a solid and the required phase for analysis is a liquid, the process is called solid-liquid extraction. A simple and broadly applicable form of solid-liquid extraction entails combining the solid with a solvent in which the analyte is soluble. Through agitation, the analyte partitions into the liquid phase, which may then be separated from the solid through filtration. The choice of solvent must be made based on the solubility of the target analyte, and on the balance of cost, safety, and environmental concerns.

 Organic Chemistry

Annexin V and Propidium Iodide Labeling

JoVE 5650

Staining with annexin V and propidium iodide (PI) provides researchers with a way to identify different types of cell death—either necrosis or apoptosis. This technique relies on two components. The first, annexin V, is a protein that binds certain phospholipids called phosphatidylserines, which normally occur only in the inner, cytoplasm-facing leaflet of a cell’s membrane, but become “flipped” to the outer leaflet during the early stages of apoptosis. The second component is the DNA-binding dye molecule PI, which can only enter cells when their membranes are ruptured—a characteristic of both necrosis and late apoptosis.

This video article begins with a review of the concepts behind annexin V and PI staining, and emphasizes how differential patterns of staining can be used to distinguish between cells progressing down different death pathways. We then review a generalized protocol for this technique, followed by a description of how researchers are currently using annexin V and PI staining to better understand cell death.

 Cell Biology

Reconstitution of Membrane Proteins

JoVE 5693

Reconstitution is the process of returning an isolated biomolecule to its original form or function. This is particularly useful for studying membrane proteins, which enable important cellular functions and affect the behavior of nearby lipids. To study the function of purified membrane proteins in situ, they must be reconstituted by integrating them into an artificial lipid membrane.

This video introduces membrane protein reconstitution concepts and related procedures, such as protein isolation using detergent, formation of artificial vesicles using lipids, incorporation of the isolated protein into the artificial vesicle, and separation of the detergent from the solution. Finally, two applications are covered: reconstitution of membrane transport proteins and reconstitution of light-harvesting proteins. 

Reconstitution is the process of restoring an isolated biomolecule to its original form or functionality. This approach is often used when studying membrane proteins, which enable many important cellular processes and affect the behavior of neighboring lipids. However, the complexity of the cell environment makes membrane protein functions difficult to study in situ. The proteins can be extracted and purified, but their actual functions cannot be evaluated without a membrane. Therefore, isolated membrane proteins are reconstituted by integration into an artificial lipid membrane, such as a liposome. This video will introduce the principles of membrane protein reconstitution, a general reconstitution procedure, and a few applications in biochemistry.

Cell membranes primarily consist of phospholipids and membrane proteins. The phospholipids form a bilayer in which the hydrophilic phosphate heads interact with the aqueous interior and exterior of the cell, while the hydrophobic fatty acid tails interact with each other in the bilayer.

Some membrane proteins only interact with the membrane by electrostatic or noncovalent interactions. Others, called 'integral proteins', are embedded in the lipid bilayer.

Like the bilayer, integral proteins have hydrophilic ends and a hydrophobic center, and are held in place by hydrophobic interactions. Integral proteins that span the entire membrane are known as 'transmembrane proteins'.

The interactions between these proteins and the membrane are so strong that even lysing the cells will not separate them. A special surfactant called a detergent is used to extract the proteins. Similar to phospholipids, detergents have hydrophilic heads and lipophilic tails, and can enter the membrane freely.

Inside the membrane, the lipophilic tails of the detergent interact with the hydrophobic protein core. This surrounds the protein with a shell of the hydrophilic detergent heads, which disrupts the protein-lipid interactions.

The protein-detergent complex is now easily separated from the membrane. The detergent makes the complex soluble in aqueous solutions, and ready for reconstitution in an artificial membrane.

Proteins are often reconstituted in the membranes of liposomes, which are artificial vesicles. To prepare liposomes, dried lipids are hydrated and agitated to induce vesicle formation. When a detergent is added, it is incorporated into the liposome membranes.

To reconstitute the protein, the solubilized proteins and liposomes are combined, and then the detergent is removed from solution by dialysis or chemical adsorption. The proteins and liposomes rapidly assemble into proteoliposomes, so only the hydrophilic groups are exposed. The proteins then function as they would in a cell membrane, and can be investigated in isolation.

Now that we've covered the basics of protein reconstitution, let's go over a protocol for reconstituting membrane proteins in liposomes.

To begin isolating the membrane proteins, the cells are lysed. Unbroken cells are removed with centrifugation.

The supernatant is centrifuged at a higher speed to pellet the membranes. The pellet is re-suspended and a detergent is added to extract the proteins.

The remaining cell debris is removed by additional centrifugation. The protein is purified from the supernatant with column chromatography and then concentrated or purified further as needed.

To begin preparing the liposomes, a suspension of phospholipids in organic solvent is dried under nitrogen or argon.

The phospholipids are hydrated with hydration buffer, and the mixture is sonicated to finish creating the liposomes.

Detergent is added to solubilize the liposomes, which is then combined with the proteins.

The detergent is then removed by adsorption onto polystyrene beads, dialysis, or a detergent-binding column. The resulting proteoliposomes are ready to be purified and used in subsequent experiments.

Now that you are familiar with the basics of a membrane protein reconstitution procedure, let's look at a few applications of protein reconstitution in biochemistry.

A membrane-transport protein was reconstituted to gain a clearer understanding of its transport mechanism. Its function post-reconstitution was verified with an efflux of iodide ions. Then, the transport activity was studied in the presence of various small molecule ion channel inhibitors and potentiators. In this way, the direct interactions of these small molecules with the transport protein could be studied.

Chlorophyll and carotenoid-binding membrane proteins in plants harvest light, promote charge separation, and mitigate light damage. By reconstituting these light-harvesting proteins, their folding dynamics and interaction with pigments can be studied. The light-harvesting proteins reconstituted with this technique had very similar optical properties to the native proteins. Fluorescence emission spectroscopy can then be used to study energy transfer from pigments to the reconstituted light-harvesting proteins.

You've just watched JoVE's video on reconstitution of membrane proteins. Reconstitution is a way to transfer important proteins to a cell mimic for further investigation. This video covered the principles of protein reconstitution, a reconstitution protocol, and a few applications in biochemistry. Thanks for watching!

 Biochemistry

Synthetic Biology

JoVE 5792

This video presents synthetic biology and its role in bioengineering. Synthetic biology refers to the methods used to genetically modify organisms in order to make them capable of producing large quantities of a product. This product could be a protein that the cell already makes or a new protein that has been encoded in a newly-inserted DNA sequence.

Here, we discuss how an organism's genetic material is modified using transformation or transfection. Then, the process is shown in the laboratory, and the applications of the technique discussed.

 Bioengineering

Energy Diagrams, Transition States, and Intermediates

JoVE 11701

Free-energy diagrams, or reaction coordinate diagrams, are graphs showing the energy changes that occur during a chemical reaction. The reaction coordinate represented on the horizontal axis shows how far the reaction has progressed structurally. Positions along the x-axis close to the reactants have structures resembling the reactants, while positions close to the products resemble the products.  Peaks on the energy diagram represent stable structures with measurable lifetimes, while other points along the graph represent unstable structures that cannot be isolated.

This high-energy unstable structure is called the transition state or activated complex. In this high-energy process, bonds are in the process of being broken and/or formed simultaneously. The structure is so strained that it transitions into new, less strained structures.

George Hammond formulated a principle that relates the nature of a transition state to its location on the reaction diagram. The Hammond Postulate states that a transition state will be structurally and energetically similar to the species nearest to it on the reaction diagram. In the case of an exothermic reaction, the transition state resembles the reactant species, whereas, in the case of an endothermic reaction, the transition state resembles the products. In a multi-step reaction, each step has a transition state and corresponding activation energy. The transition states of such reactions are punctuated with reactive intermediates, which are represented as local minima on the energy diagrams.

Reactive intermediates are products of bonds breaking and cannot be isolated for prolonged periods of time. Some of the most common reactive intermediates in organic chemistry are carbon ions or radicals. Carbocations are electrophiles, and carbanions are nucleophiles. Carbon radicals have only seven valence electrons and may be considered electron deficient; however, they do not, in general, bond to nucleophilic electron pairs, so their chemistry exhibits unique differences from that of conventional electrophiles. Radical intermediates are often called free radicals.

 Core: Organic Chemistry

Nucleophiles

JoVE 11744

The word “nucleophile” has a Greek root and translates to nucleus-loving. Nucleophiles are either negatively charged or neutral species with a pair of electrons in a high-energy occupied molecular orbital (HOMO). As these species tend to donate electron pairs, nucleophiles are considered Lewis bases as well. Negatively charged species, like OH, Cl, or HS, with one or several pairs of electrons, are typically nucleophiles. Similarly, neutral species such as ammonia, amines, water, and alcohol have non-bonding lone pairs of electrons and can act as nucleophiles. Furthermore, molecules without a lone pair of electrons can still act as nucleophiles, such as alkenes and aromatic rings with bonding π orbitals.

The relative strength of a nucleophile to displace a leaving group in a substitution reaction is called nucleophilicity. The negatively charged species are more nucleophilic than their neutral counterpart species. As an empirical rule, the higher the pKa of a conjugate acid, the better the nucleophile. For example, the hydroxide ion — a conjugate base of water (pKa 15.7) is a better nucleophile than the acetate ion — a conjugate base of acetic acid (pKa~5).

Since nucleophilicity is not the inherent property given to a specific species, it is affected by many factors, including the type of solvent in which the reaction is conducted. In polar protic solvents, high solvation of anions reduces the nucleophile’s availability to participate in substitution reactions.

When comparing halides, fluoride, being the smallest and most electronegative anion, is solvated the strongest, while iodide, the largest and least electronegative ion, is solvated the least. Thus, in polar protic solvents, iodide is the best nucleophile. In polar aprotic solvents, however, anions are “naked” due to poor solvation and can participate freely in a nucleophilic attack. In polar aprotic solvents, the nucleophile’s basicity dictates its nucleophilicity making fluoride the best nucleophile.

Furthermore, the polarizability of atoms affects nucleophilicity. Polarizability describes how easily electrons in the cloud can be distorted. A nucleophile with a large atom has greater polarizability, meaning it can donate a higher electron density to the electrophile compared to a small atom, whose electrons are held more tightly.

 Core: Organic Chemistry

Hydroboration-Oxidation of Alkenes

JoVE 11779

In addition to the oxymercuration–demercuration method, which converts the alkenes to alcohols with Markovnikov orientation, a complementary hydroboration-oxidation method yields the anti-Markovnikov product. The hydroboration reaction, discovered in 1959 by H.C. Brown, involves the addition of a B–H bond of borane to an alkene giving an organoborane intermediate. The oxidation of this intermediate with basic hydrogen peroxide forms an alcohol.

Figure1

Borane as a reagent is very reactive, as the boron atom has only six electrons in its valence shell. The unoccupied 2p orbital of the boron is perpendicular to the plane, which is occupied by the boron and the three other hydrogens oriented at an angle of 120°. Thus, borane is electrophilic with its structure resembling a carbocation without any charge.

Figure2

Due to high reactivity, two borane molecules dimerize such that two hydrogen atoms are partially bonded to two boron atoms with a total of two electrons. Hence, they are called three-center, two-electron bonds. Diborane co-exists in equilibrium with a small amount of borane. 

Figure3

The electron-deficient borane easily accepts an electron pair from tetrahydrofuran to complete its octet forming a stable borane-ether complex. This is used as a reagent in hydroboration reactions under an inert atmosphere to avoid spontaneous ignition in the air.

Figure4

Hydroboration Mechanism

The mechanism starts with a borane attacking the π bond at the less substituted and sterically less hindered site of an alkene forming a cyclic transition state. The overall result is a syn-addition of BH2 and hydrogen across the alkene double bond, producing an alkylborane. The reaction of a second alkene with the alkylborane yields a dialkylborane followed by the addition of a third alkene to produce a trialkylborane.

Figure5

Oxidation Mechanism

The oxidation begins with the deprotonation of hydrogen peroxide by a hydroxide ion forming a hydroperoxide. The hydroperoxide acts as a nucleophile and attacks the trialkylborane, resulting in an unstable intermediate. This is followed by the migration of an alkyl group from boron to the adjacent oxygen atom, releasing a hydroxide ion. This series of three steps is repeated to convert the remaining trialkylborane into a trialkoxyborane.

The boron atom of the trialkoxyborane is attacked by the nucleophilic hydroxide ion with the subsequent departure of the alkoxide ion neutralizing the formal charge on the boron atom. Finally, protonation of the alkoxide ion gives an alcohol as the final product.

Figure6

 Core: Organic Chemistry

Electrophilic Addition to Alkynes: Hydrohalogenation

JoVE 11837

Electrophilic addition of hydrogen halides, HX (X = Cl, Br or I) to alkenes forms alkyl halides as per Markovnikov's rule, where the hydrogen gets added to the less substituted carbon of the double bond. Hydrohalogenation of alkynes takes place in a similar manner, with the first addition of HX forming a vinyl halide and the second giving a geminal dihalide.

Figure1

Addition of HCl to an Alkyne

Mechanism I – Vinylic carbocation Intermediate

The mechanism begins with a proton transfer from HCl to the alkyne. Here, the π electrons attack the hydrogen atom of hydrogen chloride and displace the chloride ion. This gives a stable secondary vinylic carbocation, which is further attacked by the chloride ion to form a vinyl chloride.

Figure2

The addition of a second equivalent of hydrogen chloride proceeds with the protonation of vinyl chloride, giving two possible carbocations. In the secondary carbocation the positive charge is delocalized through resonance. As a result, it is more stable and favored over the primary carbocation.

Figure3

Figure4

A second nucleophilic attack by the chloride ion gives a geminal dichloride.

Figure5

Mechanism II – Concerted Termolecular Process

In a vinylic carbocation, the positive charge resides on an electronegative sp-hybridized carbon making it unstable. The second mechanism avoids the formation of this carbocation. Instead, it proceeds via a termolecular (three molecules) process where the alkyne interacts simultaneously with two equivalents of a hydrogen halide like HCl. This leads to a transition state with a partially broken C–C π bond and partially formed C–Cl and C–H σ bonds. The net result is a trans addition of hydrogen from one HCl and a chloride from the other HCl to form a chloroalkene. This further reacts with the displaced HCl to form a geminal dichloride as the final product.

Figure6

Halogenation of Alkynes with Peroxides

When treated specifically with HBr in the presence of peroxides, terminal alkynes undergo anti-Markovnikov's addition. The Br gets added to the less substituted carbon forming a mixture of E and Z alkenes.

Figure7

 Core: Organic Chemistry

The Structure of Intermediate Filaments

JoVE 11916

The intermediate filaments are one of three widely studied cytoskeletal filaments. They are so named as their diameter (10 nm) is in between that of microfilaments (7 nm) and the microtubules (25 nm).  These filaments are highly stable and can remain intact when exposed to high salt concentrations and detergents. These filaments are responsible for providing stability and mechanical support to the cells. They also help in cell adhesion and maintaining tissue integrity.

Intermediate filaments are found in almost all eukaryotic cells except lower eukaryotes like fungi and invertebrates like arthropods. In humans, around 70 genes code for different types of intermediate filament with cell type-specific expression depending on the function they perform. The mutation in these genes can lead to different diseases or disorders in humans, including Werner's syndrome, Alexander's disease, and Charcot-Marie tooth disease.

Intermediate filaments are non-polar with no defined plus or minus ends like microfilaments and microtubules; thus, no molecular motor proteins are known to be associated with them. Unlike microfilaments (F-actins) and microtubules made up of globular proteins, the monomeric units of intermediate filaments are rigid, fibrous rope-like proteins. These monomers vary among the cell types, forming a specific type of intermediate filament depending on their function. However, all monomeric units comprise a conserved alpha-helical central coiled-coil rod domain flanked by the head and tail domains. The conserved alpha-helical rod domain has 310 amino acids divided into three conserved segments with hydrophobic amino acid sequences separated by linkers. The rod domain helps in the lateral association of monomers to form the dimers and subsequently the tetramers, the basic soluble units of the intermediate filaments. The tetramers then assemble into unit-length filaments, which further associate to form the intermediate filaments.

 Core: Cell Biology

Structure and Nomenclature of Thiols and Sulfides

JoVE 12111

Thiols and sulfides are sulfur analogs of alcohols and ethers, respectively, where the sulfur atom takes the place of the oxygen atom. Thus, thiols are generally represented as RSH, where R is an alkyl substituent and —SH is the functional group. On the other hand, in sulfides, the central sulfur atom is bonded to two hydrocarbon groups on either side. Depending upon the type of group, sulfides can be either symmetrical or asymmetrical. Both thiols and sulfides display a bent geometry, similar to alcohols and ethers. However, the larger size of the sulfur atom deviates the C–S–C and C–S–H bond angles and the C–S and S–H bond lengths from the corresponding bond angle and bond length values of alcohols and ethers.

Generally, common names of thiols consist of the suffix "mercaptan," following the name of the parent alkyl group. In contrast, the IUPAC names of thiols are obtained by adding the suffix "thiol" to the parent alkane name while retaining the letter "e" of the parent alkane. Similarly, the suffix "sulfide" is added to the alkyl groups' names listed in alphabetical order to derive the common names of sulfides. The IUPAC names of sulfides are obtained by first determining the parent chain and prefixing the name of the smaller alkyl group as an alkylthio substituent to the parent alkane name.

 Core: Organic Chemistry

Diels–Alder vs Retro-Diels–Alder Reaction: Thermodynamic Factors

JoVE 12418

The Diels–Alder reaction is thermally reversible, meaning that the reaction reverts to the starting diene and dienophile under suitable temperatures. The forward reaction gives a cyclohexene derivative and is favored at low to medium temperatures. The reverse process, also called retro-Diels–Alder reaction, is a ring-opening process favored at high temperatures.

Figure1

Thermodynamic factors

The influence of temperature on the spontaneity of a particular reaction can be assessed based on the change in the Gibbs free energy, ΔG. If ΔG is negative, the reaction occurs spontaneously. However, if ΔG is positive, the reaction occurs spontaneously in the opposite direction.

Figure2

ΔG is a composite of two terms, ΔH and −TΔS. In a Diels–Alder reaction, two stronger σ bonds are formed at the expense of two weaker π bonds, resulting in a negative ΔH. Cycloaddition reactions proceed with a decrease in entropy. Consequently, ΔS is negative, and the term −TΔS is positive. At low temperatures, the sum of the two terms is negative, implying that the forward reaction is spontaneous. In contrast, the sum is positive at high temperatures, indicating that the reverse reaction is spontaneous.

 Core: Organic Chemistry

Lineage Commitment

JoVE 12513

Commitment is the  process whereby stem cells:

  1.  lose their ability to form all cell types and
  2. irreversibly change into a specific type.

The multipotent hematopoietic stem cells, (HSCs), differentiate into the multipotent hematopoietic progenitor cells,  (HPCs). The HPCs express many lineage-specific cytokine receptors. Each of these receptors binds specific cytokines, activates distinct signaling pathways, and expresses a particular gene set. The HPCs further differentiate to form committed progenitors, forming either  common myeloid progenitors (CMPs) or common lymphoid progenitors (CLPs). The CMPs and CLPs proliferate, self-renew, and further differentiate into mature blood cells and immune cells depending on the receptors they express and the specific cytokines that bind. For example:

  • Thrombopoietin (TpoR)- They help progenitors differentiate into megakaryocytes or platelets.
  • Erythropoietin (EpoR)-These receptors promote the development of erythrocytes.
  • Macrophage-colony stimulating factor (M-CSFR)-These receptors regulate the formation of macrophages.
  • Granulocyte-colony stimulating factor (G-CSFR)- They help HPCs differentiate into granulocytes.
  • Interleukin-7 receptor (IL-7R)- They help progenitor cells become lymphocytes.

Thus, lineage commitment helps HSCs lose their multipotency and differentiate into more restricted cell fate.

 Core: Cell Biology

Mechanism of Ciliary Motion

JoVE 13108

The ciliary structures were first seen in 1647 by Antonie Leeuwenhoek while observing the protozoans. In lower organisms, these appendages are responsible for cell movement, while in higher organisms, these appendages help in the movement of the extracellular fluids within the body cavities.

The cilia are made up of microtubules in a 9+2 arrangement, with nine microtubule doublet ring bundles, surrounding a pair of central singlet microtubule bundles. The doublet microtubule bundles are connected by nexin protein and axonemal dyneins. Radial spokes connect these outer doublet microtubules to the inner central pair. The coordinated movement of axonemal dyneins is responsible for the characteristic whip-like movement of the cilia. This characteristic ciliary movement is explained by the switch inhibition or switch-point mechanism proposed by Wais-Steider and Satir in 1979. The model suggests that during the ciliary motion, only half of the dyneins are active at a given time, while the other remains inactive. The axonemal dyneins on either side alternately switch between active and inactive forms to propel the ciliary motion. The sliding microtubules within the cilia require energy from ATP hydrolysis within the heavy chain domain of the axonemal dyneins.

Cilia in humans move rhythmically; they constantly remove waste materials such as dust, mucus, and bacteria through the airways, away from the lungs, and toward the mouth. Beating cilia on cells in the female fallopian tubes move egg cells from the ovary towards the uterus. A flagellum is an appendage larger than a cilium and specialized for cell locomotion. In humans, sperms are the only flagellated cell that must propel themselves towards female egg cells during fertilization.

 Core: Cell Biology

Immunoprecipitation

JoVE 13377

Immunoprecipitation, or IP, is a widely used technique that employs protein-antibody interactions to isolate proteins or protein complexes in their native state for studying protein-protein interactions, quaternary structures, or supramolecular complexes. Various modifications of the technique, including chromatin IP, cross-linking IP, and fluorescence IP, are commonly used.

Chromatin Immunoprecipitation

Chromatin immunoprecipitation, also known as ChIP, is used to study protein-DNA or protein-RNA interactions. It is an important technique for studying crucial cellular processes such as gene transcription, DNA replication, recombination and repair, cell cycle progression, and epigenetics. It is also helpful to identify the in vivo location of binding sites of various transcription factors, histones, and other regulatory proteins.

The major steps in ChIP include fixation, sonication, immunoprecipitation, and analysis. It involves cross-linking the target protein with the DNA using a fixing agent, such as formaldehyde, followed by sonication or enzymatic hydrolysis to obtain smaller chromatin fragments. The technique then utilizes high-affinity antibodies specific against the protein of interest to capture the DNA bound to the protein in an immunoprecipitation reaction. The cross-linking is then reversed using high heat, high salt concentration, and proteinase K to release the DNA from associated proteins. The DNA is further purified to prepare it for analysis.

Cross-linking Immunoprecipitation

Cross-linking immunoprecipitation, or CLIP, identifies the regions of protein binding sites on endogenous RNA by co-precipitating them in the active transcription phase. RNA molecules are cross-linked to proteins to hold them together tightly and prevent their degradation. The procedure for the breakdown and isolation of the complexes is similar to ChIP. This technique is used to study the interaction of RNA with RNA binding proteins and their modifications in various biological systems.

Limitations of IP

Though immunoprecipitation techniques tend to reduce the number of purification steps, it has certain limitations. The antibody binding recombinant bacterial proteins, proteins A or G, conjugated to agarose beads, and antibodies used in the method may undergo non-specific binding, introducing contaminants in the purified protein preparation. Also, the immobilization of antibodies on beads requires optimization and is time-consuming. The washing of the beads after the antigen-antibody complex is critical, but there is a chance of losing the target protein at this step.

 Core: Cell Biology

Immunocytochemistry and Immunohistochemistry

JoVE 13393

Immunocytochemistry (ICC) and immunohistochemistry (IHC) are techniques that use antibodies to check for specific proteins or antigens in a sample. The technique was first published by Albert Coons in 1941 to detect the presence of pneumococcal antigen in tissue sections from mice infected with Pneumococcus. Immunocytochemistry helps localization of proteins or antigens in individual cells like blood cells, stem cells, etc., while immunohistochemistry does the same for tissue samples.

These techniques can be colorimetric if enzyme-conjugated antibodies are used or if fluorescence-based fluorophore-tagged antibodies localize the proteins or antigens of interest. The colored product obtained via the enzymatic reaction or fluorescence emitted by the fluorophore is then visualized under an optical or immunofluorescence microscope.

These techniques can use primary or secondary antibodies tagged with an enzyme or a fluorophore. If a primary antibody is used, the technique is referred to as direct ICC or IHC, as the conjugated antibodies directly bind to the epitope of the protein or antigen of interest. In indirect, a secondary antibody tagged with the enzyme or fluorophore binds to the primary antibody attached to the epitope of the protein or antigen of interest. The signal emitted upon binding of a secondary antibody is stronger than the signal from a primary antibody.

Both the techniques vary slightly in the steps of sample preparation. In ICC, a monoculture layer of cells is generated, which are then treated with fixing agents like formaldehyde to prevent the enzymatic degradation of the antigens, followed by treatment with detergents like Triton X or Tween 20 to make the cell membrane permeable for the entry of antibodies. Similarly, in IHC, the tissue samples are first treated with a fixing solution, followed by embedding in paraffin wax to prevent damage to fragile tissue. The tissue samples are then sectioned into thin slices and treated with detergent before the introduction of antibodies.

Once the cells or tissue slides are treated with detergent, the antibodies are then added. In direct ICC or IHC, the primary antibody tagged with fluorophore or enzymes is added, and post-incubation, they are visualized under immunofluorescence or optical microscope. In indirect ICC or IHC, the samples are first incubated with the primary antibody, followed by washing to remove unbound antibodies. The secondary antibody is then introduced, which binds with the primary antibody and emits fluorescence or colored product upon adding dye. The samples are now ready for visualization.

 Core: Cell Biology

Overview of Regeneration and Repair

JoVE 13473

Regeneration and repair processes are critical in healing damages caused by injury, disease, and aging. In regeneration, the damaged tissue is entirely replaced with new growth that restores the original architecture and function. In contrast, tissue repair usually results in a fixed tissue architecture involving scar formation. Scars generally do not reestablish tissue function and may also exhibit structural abnormalities at the injury site.

Regeneration

All animals have varying degrees of regenerative capabilities, but only some animals exhibit exceptional regeneration, such as generating an entire organ. For example, salamanders have a specialized population of stem cells that plays a significant role in regenerating fully functional limbs and tails. In mammals, only a few adult tissues like the liver, epithelium of the gut, and bone marrow can regenerate after an injury. The planarian flatworms, with a much simpler body structure, are capable of whole-body regeneration; if a planarium is cut into several pieces, each piece can develop into a completely new organism. Similarly, the hydra can regenerate its whole body even from a small fragment.

The regeneration process begins with the crucial step of covering the injury site. The epidermal cells migrate to the injury site, forming a thick covering called the apical epidermal cap. This is followed by clustering of undifferentiated fibroblast cells underneath the cap, forming a blastema. The blastema is a population of stem cells that has the potential to differentiate  into any tissue or organ. The blastema is supplied with oxygen and nutrients via a new blood capillary network, allowing the cells to divide and differentiate.

Repair

The repair mechanism involves four phases: hemostasis, inflammation, proliferation, and remodeling. After an injury, the exposure of collagen— a structural protein in the walls of blood vessels— begins a coagulation cascade and vasoconstriction, thus, minimizing blood loss. During hemostasis, the clot fills the wound bed, providing a temporary matrix for the movement of various cells such as macrophages, neutrophils, and platelets. Platelet degranulation activates the complement cascade that stimulates inflammatory cells to attack the bacteria. During the proliferation phase, various cytokines and growth factors are released into the wound site, signaling the cells to proliferate and heal the wound. The final repair phase involves forming scar tissue that marks the completion of repair.

 Core: Cell Biology

Heat and Free Expansion

JoVE 13677

The work done by a thermodynamic system depends not only on the initial and final states but also on the intermediate states—that is, on the path. Like work, when heat is added to a thermodynamic system, it undergoes a change of state, and the state attained depends on the path from the initial state to the final state. Consider an ideal gas cylinder fitted with a piston. When the cylinder is heated at a constant temperature, the gas molecules absorb energy and expand slowly in a controlled isothermal manner. This pushes the piston upwards, and gas eventually attains the final volume. Here, the work is done by the gas due to heat expansion on the piston.

The gas can also attain the same final volume through a different process. Consider a cylinder surrounded by insulating walls and divided by a thin, removable partition. The lower compartment is filled with the same amount of gas at the same temperature so that the initial state is the same as mentioned above. When the partition is removed, the gas undergoes a rapid, uncontrolled expansion, with no heat passing through the insulating walls, and reaches the same final volume as in the above case. Here, no work is done by the gas during this expansion as it does not push against anything that moves. This uncontrolled expansion of a gas into the vacuum is called a free expansion.

Experimentally, there is no temperature change under a free expansion of ideal gas. This means that the final state of the gas remains the same. The intermediate states (pressures and volumes) during the transition from the initial to the final state are entirely different in the two cases. As a result, it represents two different paths connecting the same initial and final states.

 Core: Physics

Systematic Error: Methodological and Sampling Errors

JoVE 14507

In the case of systematic errors, the sources can be identified, and the errors can be subsequently minimized by addressing these sources. According to the source, systematic errors can be divided into sampling, instrumental, methodological, and personal errors.

Sampling errors originate from improper sampling methods or the wrong sample population. These errors can be minimized by refining the sampling strategy. Defective instruments or faulty calibrations are the sources of instrumental errors. Periodic calibration of instruments is essential to eliminate such errors.

Method errors occur due to the limitations in an analytical method, the non-ideal behavior of reagents used, and invalid assumptions made while setting up the measurement. These errors can be mitigated using standard reference materials or analysis independent of the existing method but carried out in parallel.

Personal errors arise due to analysts' carelessness or lack of skill. Proper organization of materials and equipment, standardization of protocols, and attention to detail can help minimize these errors. Automated procedures can also be instituted to minimize human handling, reducing personal errors.

Systematic errors can also either be constant errors or proportional errors. The absolute magnitude of constant errors remains the same irrespective of the sample size. These errors can be minimized by increasing the sample size, as the contribution from the constant error relative to a larger sample size is less significant than the same constant error relative to a smaller sample size. On the other hand, a proportional error will increase in magnitude with increasing sample size, hence the term 'proportional.' Increasing the sample size will not help to reduce these types of errors. Using high-precision instruments is one way of reducing proportional errors.

 Core: Analytical Chemistry

Ionic Strength: Overview

JoVE 14523

The ionic strength of a solution is a quantitative way of expressing the total electrolyte concentration of a solution. This concept was first introduced in 1921 by two American physical chemists, Gilbert N. Lewis and Merle Randall, while describing the activity coefficient of strong electrolytes. During the calculation of ionic strength (I or μ), all the cations and anions are considered. However, the concentration (c) of an ion with a greater charge number (z) has a greater contribution to the total ionic strength because the charge of the ion is squared.

While calculating ionic strength for a salt that will produce multiple equivalents of the same ion on dissociation, we need to account for the contribution from each ion. For example, the ionic strength of 0.1 mol/L Na2SO4 can be calculated as follows:

The concentration of Na+ is 0.2 mol/L because one molecule of Na2SO4 will dissociate to give two Na+ ions in solution. The ionic strength of dilute solutions can be calculated easily. However, in a more concentrated solution, the calculation becomes more complex and less accurate, as the salts do not dissociate completely. For example, in an aqueous solution of 0.025 mol/L MgSO4, 25% to 35% of MgSO4 exists as the ion pair MgSO4(aq).

The concept of ionic strength can be further extended to strong and weak acids. Because strong acids dissociate completely in solution, their ionic strengths can be calculated in the same way as those of dissociated salts. For weak acids, the concentration of ionized species can be calculated from the ionization constant value and then used for ionic strength determination. If the acid is very weak and mostly remains non-ionized, its contribution to the total ionic strength of the solution is negligible.

 Core: Analytical Chemistry

Titrimetric Methods: Types and Commonly Used Strategies

JoVE 14539

In chemistry, titrimetric methods are broadly classified into three types: volumetric, gravimetric, and coulometric. Volumetric titrations involve measuring the volume of a titrant of known concentration that is required to react completely with an analyte. In gravimetric titrations, the standard solution reacts with the analyte to form an insoluble precipitate, which is filtered, dried, and weighed. In coulometric titrations, current is applied to an electrochemical reaction until the reaction is complete, and the charge required for this reaction is calculated to determine the amount or concentration of the reactant(s).

Generally, in titrations, a titrant with a known concentration is added to the analyte until the analyte is consumed, i.e. when the endpoint is reached. In this case, the process is called direct titration. In some cases, the reaction may be too slow for a sharp endpoint to be obtained. In such cases, back titration is used. In a back titration, the analyte is allowed to react with an excess amount of a standard solution, and the remaining quantity of this standard solution is titrated with a second standard solution. The exact amount of the first standard solution required to react with the analyte can then be calculated.

Sometimes, no suitable indicator can be found for either a direct or an indirect (back) titration. In this case, displacement titration can be used, in which the analyte displaces a reagent, usually from a complex, and the displaced reagent is measured by titration.

 Core: Analytical Chemistry

Masking and Demasking Agents

JoVE 14578

EDTA titrations may necessitate masking and demasking agents to temporarily protect a particular metal ion in a mixture from the EDTA reaction. These agents facilitate the sequential analysis of the metal ions by forming stable complexes with some—but not all—metal ions during certain steps.

There are many masking agents, such as cyanide, fluoride, triethanolamine, thiourea, and 2,3-bis(sulfanyl)propan-1-ol (formerly 2,3-dimercapto-1-propanol), with the masking agent chosen based on the metal ions involved. For example, a cyanide mask is used during the EDTA titration of a lead and cadmium mixture. Cyanide only masks the cadmium ion, so the lead ion still reacts with EDTA.

Demasking agents are used to release the metal ions from masking agents. For example, formaldehyde acts as a demasking agent for cyanide complexes.

 Core: Analytical Chemistry

Inductively Coupled Plasma–Mass Spectrometry (ICP–MS): Overview

JoVE 14594

In inductively coupled plasma–mass spectrometry (ICP–MS), an inductively coupled plasma (ICP) torch is used as an atomizer and ionizer. Solid samples are dissolved and volatilized before being introduced into the high-temperature argon plasma, while solution samples are nebulized and passed through the high-temperature argon plasma. Plasma dissociates the analytes and ionizes their component atoms to form a mixture of positive ions and molecular species. The positive ions are then passed on to the mass spectrometer and analyzed by their m/z values. Since an ICP source operates at atmospheric pressure and a mass spectrometer requires a vacuum, the plasma is passed through an interfacial region, reducing the pressure from 1 bar to around 10−9 bar. The interfacial region consists of two metallic cones (a sampler cone and a skimmer cone), an extraction lens, and a collision cell. The sampler cone and skimmer cone have narrow orifices allowing only a small amount of plasma to pass through. The high negative voltage of the extraction lens allows only the positive ions to enter the collision cell. The collision cell reduces the range of kinetic energies of the ions and directs the ion beam into the quadrupole mass analyzer. ICP–MS is a useful technique in elemental analysis due to its low detection limits, high selectivity, and excellent sensitivity.

 Core: Analytical Chemistry

Gravimetry: Inorganic And Organic Precipitating Agents

JoVE 14731

In gravimetry, the precipitant is chosen carefully to obtain a pure solid that can be easily filtered. Common inorganic precipitants can be used to determine several cations and anions. In some cases, the formation of the same precipitate can be used to determine the cation and the anion. For example, the reaction of barium and chromate ions to give barium chromate is used to determine both barium and chromate. However, precipitates such as hydroxides, oxalates, and metal ammonium phosphates are first converted to a weighable form. Precipitation methods can also be applied to determine organic functional groups such as organic halides, carbonyl, alkoxy groups, aromatic nitro, azo, and phosphate.

Organic precipitants are usually more selective than their inorganic counterparts and yield sparingly soluble precipitates with high molecular masses. A small number of analyte ions will yield a large amount of precipitate. For example, sodium tetraphenylborate is a near-specific precipitant for potassium and ammonium ions, yielding ionic precipitates. Several organic precipitants contain multiple functional groups that can bond with the cation to generate five- or six-membered rings called chelates. Typical chelating agents include 8-hydroxyquinoline and cupferron.

 Core: Analytical Chemistry

Functions of Smooth Muscles

JoVE 14856

Smooth muscles are an important type of muscle tissue that plays a vital role in the involuntary movements of internal organs. For example, they help regulate the movement of food through the gut and the flow of blood through the circulatory system.

Function of visceral smooth muscles

Visceral smooth muscle is found in the walls of all hollow organs, except the heart, and is a key player in the involuntary movements that drive the functioning of these internal organs. This tissue is arranged in layers of tightly packed fibers that communicate via gap junctions, enabling them to contract synchronously as a single unit. For example, smooth muscle in the intestine is organized into two primary layers: the longitudinal and circular layers. The longitudinal layers run along the length of the intestine, and their contraction shortens the overall length of the intestinal tract. Meanwhile, the circular layers wrap around the circumference and facilitate constriction and dilation of the inner cavity. The alternating contraction and relaxation of these layers result in peristalsis that helps propel the contents and ensures thorough mixing and breakdown, optimizing nutrient absorption and digestion.

Function of multi-unit smooth muscles

While many organs rely on the coordinated activity of visceral smooth muscle, certain areas, like the eye and walls of large arteries, contain multi-unit smooth muscle fibers. These fibers operate independently, each with its own nerve supply, allowing for precise and localized control. In large pulmonary arteries, for instance, multi-unit smooth muscles adjust the diameter of the lumen to direct blood flow to well-ventilated regions of the lung, enhancing oxygenation and adjusting to the body's changing needs.

Overall, smooth muscles in various organs demonstrate incredible adaptability and precision in their functions. For example, in the respiratory system, smooth muscles adjust the airway diameter, increasing or decreasing airflow in response to the body's oxygen demands. In the vascular system, these muscles regulate blood pressure and flow by altering the diameter of blood vessels.

 Core: Anatomy and Physiology

Muscles that Move the Forearm

JoVE 14877

The muscles that move the forearms can be divided into four groups: forearm flexors, forearm extensors, forearm pronators, and forearm supinators. The flexors and extensors act on the elbow joint, while the pronators and supinators act on the radioulnar joints.

Forearm Flexors

The biceps brachii, brachialis, and brachioradialis are forearm flexors. The biceps brachii is made up of two heads. Its long head originates at the supraglenoid tubercle of the scapula, whereas that of the short head is at the coracoid process of the scapula. The heads merge and insert at radial tuberosity. A fibrous membrane known as bicipital aponeurosis emerges from the distal end of the biceps brachii and inserts into the deep fascia of the forearm. The anterior surface of the humerus is the origin site for the brachialis muscle. It inserts into the ulnar tuberosity and coronoid process of the ulna. The brachioradialis muscle originates on the lateral border of the distal end of the humerus and inserts on the radius, superior to the styloid process.

Forearm Extensors

The triceps brachii, and anconeus are forearm extensors. The triceps brachii is a three-headed muscle. Its long head originates on the infraglenoid tubercle. The lateral head originates from the lateral and posterior surface of the humerus. The posterior surface of the humerus is also the origin point for the medial head. Together, the three heads are inserted at the olecranon of the ulna. Finally, the anconeus originates from the humerus at the lateral epicondyle and inserts on the olecranon and superior portion of the ulnar shaft.

Pronators and Supinators

Pronator teres and pronator quadratus are forearm pronators. The pronator teres originate from the medial epicondyle of the humerus and the coronoid process of the ulna. It inserts on the mid-lateral surface of the radius. The pronator quadratus originates on the distal shaft of the ulna and inserts onto the distal part of the radius. The supinator muscle supinates the forearm. It originates from the lateral epicondyle of the humerus and the ridge near the radial notch of the ulna. It inserts on the lateral surface of the proximal one-third of the radius.

 Core: Anatomy and Physiology

Propagation of Action Potentials

JoVE 14894

The propagation of an action potential refers to the process by which a nerve impulse, or "action potential," travels along a neuron.

Neurons (nerve cells) have a resting membrane potential, with a slightly negative charge inside compared to outside. This is maintained by ion channels, such as sodium (Na+) and potassium (K+) channels, which control the flow of ions. When a stimulus, like a touch or a signal from another neuron, triggers the neuron, sodium channels open, allowing sodium ions to rush into the neuron, causing depolarization.

If the depolarization is strong enough and reaches a certain threshold, it triggers an action potential. The initiation of an action potential occurs at the axon's beginning, or the initial segment, where a high concentration of voltage-gated Na+ channels allows a swift depolarization. As the depolarization advances along the axon, more Na+ channels open, facilitating the spread of the action potential. This is achieved as Na+ ions flow inwards, progressively depolarizing the cell membrane.

However, the Na+ channels become inactivated at peak depolarization, rendering them unopenable for a brief period, known as the absolute refractory period. As a result, any depolarization attempting to reverse direction is null, ensuring that the action potential's propagation is towards the axon terminals, thereby preserving neuronal polarity.

This propagation method applies to unmyelinated axons. In myelinated axons, the process differs. The depolarization spreads optimally due to the absence of constant Na+ channel opening along the axon segment. The precise placement of nodes ensures the membrane remains sufficiently depolarized at the next node.

Propagation in unmyelinated axons, known as continuous conduction, is slower due to the constant influx of Na+. In contrast, myelinated axons exhibit saltatory conduction - a faster method as the action potential leaps node to node, renewing the depolarized membrane. Furthermore, the speed of conduction can be influenced by the axon's diameter, a concept known as resistance.

 Core: Anatomy and Physiology

Diencephalon: Thalamus and Information Relay

JoVE 14910

The thalamus, often called “the gateway to the cerebral cortex,” is vital in processing and directing sensory and motor signals throughout the brain. Almost all inputs destined for the cerebral cortex, except for olfactory signals, are relayed through the thalamus. The thalamus is  a sophisticated relay station, channeling information from various brain regions to the cerebral cortex, as well as a filter, prioritizing certain signals over others based on current physiological states or needs.

The thalamus comprises several groups of paired nuclei, each having a specific role in relaying information, ensuring that signals reach their intended destinations accurately and efficiently. The paired nuclei include the following:

Anterior Nucleus: This nucleus serves as a conduit for information from the hypothalamus to the limbic system, a collection of structures involved in emotion, memory, and arousal. This pathway is essential for integrating emotional content with sensory experiences.

Medial Nuclei: Information received from the limbic system is further processed and relayed by the medial nuclei, facilitating the emotional and cognitive responses to sensory inputs.

Lateral Geniculate Nucleus: Specialized in relaying visual information, the lateral geniculate nucleus directs signals from the eyes to the visual cortex. This process is fundamental for the perception and interpretation of visual stimuli.

Medial Geniculate Nucleus: Comparable to its lateral counterpart, the medial geniculate nucleus focuses on auditory information, channeling sound signals from the ear to the auditory cortex. This relay is crucial for hearing and the cognitive processing of sounds.

Ventral Posterior Nucleus: This nucleus conveys somatosensory information, such as touch and pain, from the body to the cerebral cortex. It ensures that sensory experiences are accurately perceived and interpreted.

Lateral Dorsal Nucleus: The lateral dorsal nucleus plays a role in expressing emotions and facilitates the integration of emotional states with sensory experiences.

Lateral Posterior and Pulvinar Nuclei: These nuclei integrate sensory information, helping to combine different sensory inputs for a cohesive perception of the environment.

Beyond sensory information, the thalamus is instrumental in relaying motor impulses. Signals from the basal nuclei, which are involved in controlling movement, are relayed through the ventral anterior nuclei. Additionally, impulses originating from both the basal nuclei and the cerebellum, a region responsible for coordination and balance, are directed through the ventral lateral nuclei. This coordination ensures that motor functions are executed smoothly and accurately. The thalamus ensures the brain responds appropriately to internal and external stimuli by directing sensory and motor signals to the appropriate cortical areas.

 Core: Anatomy and Physiology

Spinal Nerves: Plexus II

JoVE 14926

The plexuses of the lower body include the lumbar, sacral, and coccygeal plexuses, which innervate the abdomen, pelvis, legs, and coccygeal region. These plexuses control the transmission of sensory information and coordinate motor functions of the lower body.

The Lumbar Plexus

The lumbar plexus is situated within the lumbar region of the back and is primarily formed by the first four lumbar spinal nerves (L1 to L4). This plexus extends its branches into several nerves, including the iliohypogastric, ilioinguinal, genitofemoral, lateral femoral cutaneous, femoral, and obturator nerves. These nerves provide sensory and motor innervation to the muscles and skin of the abdomen, pelvis, and parts of the lower limbs. The lumbar plexus facilitates movements and transmits sensory information from these areas, enabling functions such as walking, posture maintenance, and the protection of abdominal organs.

The Sacral Plexus

Within the pelvis, the sacral plexus is formed by the nerves spanning from the fourth lumbar (L4) to the fourth sacral (S4) segments. This plexus gives rise to several vital nerves, including the pudendal, superior and inferior gluteal, posterior femoral cutaneous, and sciatic nerves. Notably, the sciatic nerve is recognized as the thickest and longest nerve in the body, extending down to the lower leg and foot. This nerve is a common source of pain and discomfort when irritated or compressed — a condition called sciatica. The sacral plexus primarily provides sensory and motor innervation to the pelvis and legs, overseeing functions such as leg movement, pelvic organ control, and sensation in the lower extremities.

Due to the substantial overlap in the nerves originating from the lumbar and sacral plexuses, these two networks are sometimes collectively referred to as the lumbosacral plexus. This combined structure highlights the complexity of neural pathways and their essential roles in bodily operations.

The Coccygeal Plexus

Though the smallest among the three, the coccygeal plexus plays a specific role in innervating the skin over the coccygeal area. It is located at the very base of the spine around the coccyx and primarily comprises the anococcygeal nerve. Despite its limited size and scope compared to the lumbar and sacral plexuses, the coccygeal plexus is crucial for transmitting sensory information from the coccygeal region, contributing to the overall sensory map of the lower body.

 Core: Anatomy and Physiology

Higher Mental Functions of the Brain: Language

JoVE 14945

Language is a system of communication that allows the expression of thoughts, ideas, and feelings. The brain processes language in both hemispheres.

Language formation and comprehension take place in the dominant hemisphere. The dominant hemisphere is responsible for understanding the meaning of spoken, written, or sign language, as well as the ability to communicate. For most people, the left hemisphere is the dominant one. The right hemisphere, then, gives tone and emotional context to the language.

Language processing in the left hemisphere happens in two areas. Wernicke's area, also known as the "language comprehension center," is located in the brain's left temporal lobe and is responsible for understanding language. This area interprets incoming speech from verbal, gestures, and written sources. Broca's area, also known as the "speech production center," is located in the brain's left frontal lobe and is responsible for producing language. It sends signals to the motor cortex that then initiates movements of muscles involved in speech.

The neural pathways involved in perceiving and vocalizing a specific word operate as follows:

  • • First, information about the word is conveyed to Wernicke's area.
    •  If the word is written, Wernicke's area receives input about the word from the primary visual area.
    • If the word is spoken, Wernicke's area receives input about the word from the primary auditory area.
  • • Upon receiving the information, Wernicke's area translates the written or spoken word into the corresponding thought.
  • • Next, for the individual to respond to the spoken or written word or a gesture, Wernicke's area transmits information about the word to Broca's area.
  • • Broca's area then receives this input and develops a motor pattern to activate the necessary muscles for articulating the word.
  • • The motor pattern is conveyed from Broca's area to the primary motor area, which triggers the appropriate muscle contractions required for speech.
  • • Ultimately, the coordinated contraction of the speech muscles allows the word to be spoken.

 Core: Anatomy and Physiology

Photoreceptors and Visual Pathways

JoVE 14964

At the molecular level, visual signals trigger transformations in photopigment molecules, resulting in changes in the photoreceptor cell's membrane potential. The photon's energy level is denoted by its wavelength, with each specific wavelength of visible light associated with a distinct color. The spectral range of visible light, classified as electromagnetic radiation, spans from 380 to 720 nm. Electromagnetic radiation wavelengths exceeding 720 nm fall under the infrared category, whereas those below 380 nm are classified as ultraviolet radiation. Blue light corresponds to a wavelength of 380 nm, while dark red light corresponds to a wavelength of 720 nm. Other colors lie at varying points within this wavelength spectrum, from red to blue.

Opsin pigments, in fact, are transmembrane proteins integrated with a cofactor named retinal. This retinal is a constituent of vitamin A and a hydrocarbon molecule. The significant biochemical alteration in the extensive hydrocarbon chain of the retinal molecule is triggered when a photon impacts it. This specific process, known as photoisomerization, transitions some of the double-bonded carbons inside the chain from a cis to a trans configuration owing to the photon interaction. Before the photon interaction, the flexible double-bonded carbons of the retinal are in the cis conformation, leading to the formation of a molecule known as 11-cis-retinal. The double-bonded carbons assume the trans-conformation when a photon impacts the molecule, forming an all-trans-retinal characterized by a straight hydrocarbon chain.

The visual transduction process within the retina commences with the alteration in the retinal structure in photoreceptors. This leads to the activation of retinal and opsin proteins, which stimulate a G protein. The activated G protein then modifies the photoreceptor cell's membrane potential, causing a decrease in the release of neurotransmitters into the retina's outer synaptic layer. This state continues until the retinal molecule reverts to its original shape, the 11-cis-retinal form - a process referred to as bleaching. If a substantial amount of photopigments undergo bleaching, the retina transmits data as if contrasting visual inputs are being received. Afterimages, typically observed as negative-type images, are a common occurrence following exposure to an intense flash of light. A series of enzymatic alterations facilitate the photoisomerization reversal process, thus enabling the reactivation of the retinal in response to additional light energy.

Opsins exhibit specific sensitivity to particular light wavelengths. The rod photopigment, rhodopsin, exhibits peak sensitivity to light that has a wavelength of 498 nm. On the other hand, three color opsins are optimally responsive to wavelengths of 564 nm, 534 nm, and 420 nm, which approximately align with the primary colors—red, green, and blue. Rhodopsin found in rods demonstrates a higher sensitivity to light than cone opsins; this means that rods contribute to vision under dim light conditions while cones contribute under brighter conditions. In normal sunlight, rhodopsin is continuously bleached, and cones remain active. Conversely, in a dimly lit room, the light intensity is insufficient to stimulate cone opsins, making vision entirely reliant on rods. In fact, rods have such a high sensitivity to light that a solitary photon can trigger an action potential in a rod's corresponding RGC.

Cone opsins, differentiated by their sensitivity to distinct light wavelengths, furnish the ability to perceive color. By analyzing the responses of the three unique cone types, our brain distills color data from what we see. Consider, for instance, a bright blue light with a wavelength near 450 nm. This would cause minimal stimulation of the "red" cones, slight activation of the "green" cones, and significant stimulation of the "blue" cones. The brain computes this differential activation of the cones and interprets the color as blue. However, under dim light conditions, cones are ineffectual, and rods, which cannot discern color, dominate. As a result, our vision in low light is essentially monochromatic, meaning everything appears in varying shades of gray in a dark room.

Some common eye disorders:

Color blindness, clinically known as achromatopsia, is a condition characterized by a deficiency in distinguishing colors. This disorder usually results from an inherited defect in the retina's cones (light-sensitive cells). Symptoms may include difficulty distinguishing between colors or shades of colors.

Night blindness, medically referred to as nyctalopia or hemeralopia, is a disorder that affects an individual's ability to see in low light or at night. Causes can range from vitamin A deficiency to underlying diseases such as retinitis pigmentosa. Individuals with this disorder experience difficulties with night-time vision or adjusting to dim lighting.

Cataracts, a common eye disorder especially among older adults, are characterized by clouding of the normally transparent eye lens. This can result in blurred vision, similar to looking through a fogged-up window. Most cataracts develop slowly over time and can eventually interfere with vision.

Glaucoma is another severe eye condition where the optic nerve, which sends images to the brain, gets damaged due to increased pressure in the eye. It can lead to vision loss if left untreated. The most common type of glaucoma, open-angle glaucoma, often has no symptoms other than gradual vision loss.

 Core: Anatomy and Physiology

Hormones of the Pituitary Gland

JoVE 14980

The small, pea-sized pituitary gland is located at the base of the brain. It is crucial in regulating various bodily functions, from growth to reproduction. The gland is divided into the anterior lobe and the posterior lobe. The secretory cell clusters in the pars distalis of the anterior pituitary lobe are controlled by hypothalamic regulators and synthesize six primary hormones.

The most abundantly secreted hormone from the anterior lobe is the growth hormone, which controls overall growth by regulating the production of insulin-like growth factors from the liver, bones, and muscles. Also, it plays a vital role in regulating metabolism by balancing fat and glucose usage for energy generation.

Another group of hormones, tropins, are secreted by the anterior lobe, which controls the secretion of hormones by other endocrine glands. For example, the thyroid-stimulating hormone is a tropin regulating the hormones of the thyroid gland, while the adrenocorticotropic hormone controls cortisol release from the adrenal glands.

The luteinizing hormone, LH, and follicle-stimulating hormone, FSH, in both sexes control the functioning of reproductive organs. These hormones regulate the menstrual cycle in females and testosterone production in males. Lastly, prolactin stimulates breast milk production in lactating mothers.

 Core: Anatomy and Physiology

Power and Energy

JoVE 15065

The power and energy delivered to an element are subjects of great significance in the field of electrical engineering. It is a well-known fact that a 100-watt light bulb emits more light than a 60-watt one. Therefore, power and energy calculations play a crucial role in the analysis of electrical circuits.

Power, defined as the time rate of expending or absorbing energy, is quantified in units called watts (W). The relation between power and energy is mathematically given as

Equation1

where "p" represents the power in watts (W), "w" denotes the energy in joules (J), and "t" signifies the time in seconds (s).

The power associated with the current passing through an element within a circuit is the product of the voltage across that element and the current flowing through it. When the current enters the circuit element at the positive terminal of the voltage and exits at the negative terminal, it adheres to the passive convention. In this convention, the voltage propels a positive charge in the same direction as indicated by the current. Consequently, the power computed by multiplying the element voltage by the element current represents the power received by the element. This power is occasionally referred to as "the power absorbed by the element" or "the power dissipated by the element." Importantly, the power received by an element can assume either positive or negative values, contingent upon the actual values of the element's voltage and current.

On the other hand, when the passive convention is not adhered to, the current enters the circuit element at the negative terminal of the voltage and exits at the positive terminal. In this case, the voltage drives a positive charge in the direction opposite to that indicated by the current. Therefore, when the element voltage and current do not comply with the passive convention, the power computed by multiplying these values represents the power supplied by the element. Similar to received power, supplied power can also be positive or negative, depending on the specific values of the element's voltage and current.

The relationship between the power received by an element and the power supplied by that same element is expressed as

Equation2

When the element voltage and current align with the passive convention, the energy received by an element is given by

Equation3

Energy is defined as the capacity to perform work and is expressed in joules (J).

 Core: Electrical Engineering

Integrator and Differentiator

JoVE 15081

Op-amp circuits have significant applications in various fields, including automotive engineering. One such application is cruise control systems in cars, where op-amp circuits are integral for maintaining a constant speed. In these systems, op-amps function as both integrators and differentiators.

An integrator within an op-amp circuit produces an output directly proportional to the integral of the input signal. This is achieved by replacing the feedback resistor in a typical inverting amplifier circuit with a capacitor, resulting in an ideal integrator. An equation relating output and input voltages is derived by applying Kirchhoff's current law and utilizing current-voltage relationships for resistors and capacitors. When integrated, this equation demonstrates that the output voltage corresponds to the integral of the input signal.

Equation1

Conversely, a differentiator within an op-amp circuit yields an output proportional to the input signal's rate of change. Achieving this involves replacing the input resistor with a capacitor in a standard inverting amplifier, creating a differentiator circuit. An equation linking output and input voltages is established by applying Kirchhoff's current law and employing current-voltage relations. In this case, the equation indicates that the output voltage is proportional to the derivative of the input signal.

Equation2

It is worth noting that these op-amp circuits are valuable in energy storage applications and are often designed using resistors and capacitors due to their compactness and cost-effectiveness. While integrators are widely employed in analog computers and various applications, differentiators are less common in practice due to their tendency to amplify electrical noise, making them electronically unstable.

 Core: Electrical Engineering

Series RLC Circuit without Source

JoVE 15098

Within the field of electrical circuits, source-free RLC circuits present an intriguing domain. These circuits comprise a series arrangement of a resistor, inductor, and capacitor, operating independently of external energy sources. Their initiation hinges upon utilizing the initial energy stored within the capacitor and inductor to instigate their functionality. Their mathematical equation, a second-order differential equation, sets these circuits apart. This equation captures how the circuit's components interact, forming the basis for understanding its behavior.

Equation1

The resistor in this circuit plays a significant role by dissipating energy, leading to an exponential solution for the differential equation. Substituting this solution yields a quadratic equation, and the two roots of this equation hold special significance. These roots are the circuit's natural frequencies and are instrumental in describing its natural response.

Equation2

Equation3

Expressed in terms of the damping factor and resonant frequency, these roots provide insights into the circuit's behavior. If the damping factor surpasses the resonant frequency, the circuit exhibits an overdamped response with distinct real roots. When the damping factor equals the resonant frequency, a critically damped response ensues, characterized by equal roots. Finally, if the damping factor falls short of the resonant frequency, the circuit enters an underdamped state with complex roots.

Various response scenarios within source-free RLC circuits offer an intriguing and valuable aspect of circuit analysis. Further exploration of each case provides a comprehensive understanding of their behavior and practical applications in electrical circuits.

 Core: Electrical Engineering

Superposition Theorem for AC Circuits

JoVE 15115

Consider encountering a circuit in a steady state where all its inputs are sinusoidal, yet they do not all possess the same frequency. Such a circuit is not classified as an alternating current (AC) circuit, and consequently, its currents and voltages will not exhibit sinusoidal behavior. However, this circuit can be analyzed using the principle of superposition.

The principle of superposition stipulates that the output of a linear circuit with several concurrent inputs is equivalent to the cumulative outputs when each input operates independently. The inputs to the circuit are the voltages from the independent voltage sources and the currents from the independent current sources.

When all inputs except one are set to zero, the remaining inputs become 0-V voltage sources and 0-A current sources. Given that 0-V voltage sources equate to short circuits and 0-A current sources correspond to open circuits, the sources linked to the other inputs are replaced by open or short circuits. What remains is a steady-state circuit with a single sinusoidal input, which qualifies as an AC circuit and is analyzed using phasors and impedances.

Hence, the principle of superposition is employed to transform a circuit with multiple sinusoidal inputs at varying frequencies into several separate circuits, each with a singular sinusoidal input. Each of these AC circuits is then analyzed using phasors and impedances to determine its sinusoidal output. The aggregate of these sinusoidal outputs will coincide with the output of the initial circuit.

 Core: Electrical Engineering

Titration of Polyprotic Acids with a Strong Base

JoVE 17363

Titration of a polyprotic acid, which contains multiple ionizable protons, involves distinct dissociation steps, each with its own dissociation constant (Ka). Each successive Ka is weaker than the previous one. In the titration of a polyprotic acid like sulfurous acid with a strong base such as sodium hydroxide, the base first neutralizes the initial ionizable proton, forming an intermediate species (e.g., hydrogen sulfite ions). This step's titration curve resembles that of a weak monoprotic acid, with a half-equivalence point where the pH equals the first pKa. As the titration continues, an additional base neutralizes the second ionizable proton. This process requires twice the amount of base for complete neutralization, leading to another half-equivalence point and an equivalence point in the basic region. This pattern extends to other polyprotic acids, such as triprotic phosphoric acid, which will have three equivalence points. The number of equivalence points in the titration curve of a weak polyprotic acid corresponds to the number of ionizable protons, provided there's a significant difference between the Ka values of these protons.

 Core: Analytical Chemistry

Yeast Reproduction

JoVE 5097

Saccharomyces cerevisiae is a species of yeast that is an extremely valuable model organism. Importantly, S. cerevisiae is a unicellular eukaryote that undergoes many of the same biological processes as humans. This video provides an introduction to the yeast cell cycle, and explains how S. cerevisiae reproduces both asexually and sexually Yeast reproduce asexually through a process known as budding. In contrast, yeast sometimes participate in sexual reproduction, which is important because it introduces genetic variation to a population. During environmentally stressful conditions, S. cerevisiae will undergo meiosis and form haploid spores that are released when environmental conditions improve. During sexual reproduction, these haploid spores fuse, ultimately forming a diploid zygote. In the lab, yeast can be genetically manipulated to further understand the genetic regulation of the cell cycle, reproduction, aging, and development. Therefore, scientists study the reproduction of yeast to gain insight into processes that are important in human biology.

 Biology I

Development and Reproduction of the Laboratory Mouse

JoVE 5159

Successful breeding of the laboratory mouse (Mus musculus) is critical to the establishment and maintenance of a productive animal colony. Additionally, mouse embryos are frequently studied to answer questions about developmental processes. A wide variety of genetic tools now exist for regulating gene expression during mouse embryonic and postnatal development, which can help scientists to understand more about heritable diseases affecting human development.

This video provides an introduction to the reproduction and development of mice. In addition to clarifying the terminology used to describe developmental progression, the presentation reviews key stages of the mouse life cycle. First, major development events that take place in utero are described, with special attention given to the unique layout of early rodent embryos. Next, husbandry protocols are provided for postnatal mice, or pups, including the process of weaning, or removal of pups from their mother's cage. Since males and females must be separated at this stage to prevent unscheduled mating, the demonstration also reveals how to determine mouse sex. Subsequently, instructions are given for carrying out controlled mouse breeding, including screening for the copulatory plug, which is useful for precisely timed embryonic development. Finally, the video highlights strategies used to investigate the complex processes that govern mouse development, including the generation of genetically altered “knockout” mice.

 Biology II

fMRI: Functional Magnetic Resonance Imaging

JoVE 5212

Functional magnetic resonance imaging (fMRI) is a non-invasive neuroimaging technique used to investigate human brain function and cognition in both healthy individuals and populations with abnormal brain states. Functional MRI utilizes a magnetic resonance signal to detect changes in blood flow that are coupled to neuronal activation when a specific task is performed. This is possible because hemoglobin within the blood has different magnetic properties depending on whether or not it is bound to oxygen. When a certain task is performed, there is an influx of oxygenated blood to brain regions responsible for that function, and this influx can then be detected with specific MRI scan parameters. This phenomenon is termed the blood oxygen level ependent (BOLD) effect, and can be used to create maps of brain activity.

This video begins with a brief overview of how MRI and fMRI signal is obtained. Then, basic experimental design is reviewed, which involves first setting up a stimulus presentation that is specifically designed to test the function that will be mapped. Next, key steps involved in performing the fMRI scan are introduced, including subject safety and setting up at the scanner. Commonly used steps for data processing are then presented, including pre-processing and statistical analysis with the general linear model. Finally, some specific applications of fMRI are reviewed, such as investigations into abnormal function in psychological disorders, and combining fMRI with complimentary imaging modalities, such as diffusion tensor imaging (DTI).

 Neuroscience

An Introduction to Organogenesis

JoVE 5334

Organogenesis is the process by which organs arise from one of three germ layers during the later stages of embryonic development. Researchers studying organogenesis want to better understand the genetic programs, cell-cell interactions, and mechanical forces involved in this process. Ultimately, scientists hope to use this knowledge to create therapies and artificial organs that will help treat human diseases.

This video offers a comprehensive overview of organogenesis, including historical highlights starting with breakthrough studies done in the 1800s. Next, key questions asked by developmental biologists are introduced, followed by a discussion of how tissue transplantations, imaging, and in vitro culture techniques can be used to answer these queries. Finally, we describe how these methods are currently being employed in developmental biology laboratories.

 Developmental Biology

Assessing Dexterity with Reaching Tasks

JoVE 5424

Reaching tasks are employed in behavioral neuroscience to investigate motor learning and forelimb dexterity. Much like human hands, rodents have dexterous forepaws, which are necessary for executing coordinated and precise motor movements. Experimenters may utilize food rewards to train rodents to reach and for testing their reaching abilities. These tasks help behavioral neuroscientist in understanding how CNS injuries, such as a stroke, may impair reaching ability and dexterity in humans.

This video begins by discussing the principles and neurobiology of forelimb use in rodents, and then explains a protocol on how to conduct reaching experiments using different types of food rewards. Applications section reviews studies that involve reaching and food handling in animal models of CNS injury.

 Behavioral Science

An Overview of Genetic Analysis

JoVE 5540

An organism’s physical traits, or phenotype, are a product of its genotype, which is the combination of alleles (gene variants) inherited from its parents. To varying degrees, genes interact with each other and environmental factors to generate traits. The distribution of alleles and traits within a population is influenced by a number of factors, including natural selection, migration, and random genetic drift.

In this video, JoVE introduces some of the foundational discoveries in genetics, from Gregor Mendel’s elucidation of the genetic basis of inheritance, to how natural processes affect allele distributions within populations, to the modern synthesis of biology that brought together Mendelian genetics and Darwinian evolution. We then review the questions asked by geneticists today regarding how genes influence traits, and some of the main tools used to answer these questions. Finally, several applications of techniques such as genetic crosses, screens and evolution experiments will be presented.

 Genetics

Separation of Mixtures via Precipitation

JoVE 5558

Source: Laboratory of Dr. Ana J. García-Sáez — University of Tübingen

Most samples of interest are mixtures of many different components. Sample preparation, a key step in the analytical process, removes interferences that may affect the analysis. As such, developing separation techniques is an important endeavor not just in academia, but also in industry. 

One way to separate mixtures is to use their solubility properties. In this short paper, we will deal with aqueous solutions. The solubility of a compound of interest depends on (1) ionic strength of solution, (2) pH, and (3) temperature. By manipulating with these three factors, a condition in which the compound is insoluble can be used to remove the compound of interest from the rest of the sample.1

 Organic Chemistry

The TUNEL Assay

JoVE 5651

One of the hallmarks of apoptosis is the nuclear DNA fragmentation by nucleases. These enzymes are activated by caspases, the family of proteins that execute the cell death program. TUNEL assay is a method that takes advantage of this feature to detect apoptotic cells. In this assay, an enzyme called terminal deoxynucleotidyl transferase catalyzes the addition of dUTP nucleotides to the free 3’ ends of fragmented DNA. By using dUTPs that are labeled with chemical tags that can produce fluorescence or color, apoptotic cells can be specifically identified.

JoVE’s video on the TUNEL assay begins by discussing how this technique can be used to detect apoptotic cells. We then go through a general protocol for performing TUNEL assays on tissue sections and visualizing the results using fluorescence microscopy. Finally, several applications of the assay to current research will be covered.

 Cell Biology

Electrophoretic Mobility Shift Assay (EMSA)

JoVE 5694

The electrophoretic mobility shift assay (EMSA) is a biochemical procedure used to elucidate binding between proteins and nucleic acids. In this assay a radiolabeled nucleic acid and test protein are mixed. Binding is determined via gel electrophoresis which separates components based on mass, charge, and conformation.

This video shows the concepts of EMSA and a general procedure, including gel and protein preparation, binding, electrophoresis, and detection. Applications covered in this video include the analysis of chromatin-remodeling enzymes, a modified EMSA that incorporates biontinylation, and the study of binding sites of bacterial response regulators.

EMSA, the electrophoretic mobility shift assay, also known as the gel shift assay, is a versatile and sensitive biochemical procedure. EMSA elucidates binding between proteins and nucleic acids by detecting a shift in bands in gel electrophoresis.

This video describes the principles of EMSA, provides a general procedure, and discusses some applications.

DNA replication, transcription, and repair, as well as RNA processing are all critical biochemical processes. They all involve binding between proteins and nucleic acids. Many serious diseases and disorders are associated with modifications in this binding. EMSA is a technique for qualitatively determining whether a specific protein binds to a specific nucleic acid. First, the nucleic acid is labeled, usually with radioactive phosphorus-32, to create a probe. Then the test protein and nucleic acid probe are mixed. When a protein binds to a nucleic acid probe, the resulting complex has greater mass and a different conformation than the nucleic acid alone.

Once bound, the complexes are analyzed with gel electrophoresis. In this technique, an electric field forces macromolecules to migrate through a gel matrix. The components separate based on mass, charge, and conformation. Electrophoresis can separate protein-DNA complexes from unbound probes. Since they have different masses and conformations, they will migrate through the gel at different rates and separate. The separation is easily detected, thanks to the presence of the radioactive phosphorus, and proves the protein successfully binds to the given nucleic acid. To verify the identification of the protein, a "supershift assay" uses an antibody with a known affinity to the protein. This has the added benefit of further shifting the complex, increasing resolution.

Now that we've seen the principles, let's see it in the lab.

To begin the procedure, the protein must be isolated. To do this, molecular biology techniques are used to express the protein in cells, and then purify.

The nucleic acid is amplified and labeled to create a probe. Labeling is done through incubation for 10 min with dCTP containing radioactive phosphorus-32. A radiation-safe workbench and protective equipment are required.

The gel is then prepared. The gel needs to be non-denaturing, to prevent the protein from altering conformation and potentially unbinding from the probe during electrophoresis. Polyacrylamide gels have pore sizes of 5 to 20 nm and are useful for short probes up to 100 base pairs in length. Agarose gels have pore sizes of 70-700 nm and are useful for larger probes.

With the protein and nucleic acid probe now prepared, we proceed to binding. A TRIS buffer solution is prepared and the protein and probe are added. The pH should be similar to physiological conditions, and salt concentration sufficient to prevent the protein from forming weak bonds with non-target nucleic acids. The reaction proceeds for 20-30 min at 4 °C.

The next step is electrophoresis. A buffer of low ionic strength and a pH similar to that used in the binding reaction is utilized. It yields a "caging effect" that stabilizes complexes, increases mobility, and reduces heat generation. After electrophoresis, the components of the gel are transferred onto filter paper. In a dark room, the filter paper is then exposed to film. If the protein binds to the probe, two distinct labeled regions will be visible in the transfer. One representing the complex, and a separate one representing the unbound probe. The separation demonstrates that the protein successfully bound to the nucleic acid.

Now that we've seen the basic procedure, let's look at some applications.

Chromatin is the tightly packed complex of DNA and proteins found eukaryotic cells. Chromatin-remodeling enzymes modify the structure to open the DNA to transcription. As this changes the mobility of the complex, EMSA can be used to explore the binding activity of the enzyme.

An alternate approach to labeling takes advantage of the interaction between DNA and the methyltransferase enzyme. A cofactor can be modified to bind permanently to DNA via methyltransferase. Rather than labeling with phosphorus-32, the cofactor is conjugated to biotin, which is advantageous because it's not radioactive. Because the cofactor is site-specific in its attachment, it’s relevant to genotyping, methylation detection, and gene delivery. The biotinilated nucleic acids are detected through ultraviolet fluorescence.

When environmental stimuli activate histidine kinases, a “response regulator” is phosphorylated. This in turn binds to DNA, affecting transcription, which can be studied by EMSA. For instance, a response regulator in Desulfovibrio vulgaris was demonstrated to bind to the gene of interest. EMSA was used to verify that the binding took place.

You've just watched JoVE's video on the electrophoretic mobility shift assay. You should now understand its principles of operation, the steps in its procedure, and its major operating parameters.

Thanks for watching!

 Biochemistry

Batch and Continuous Bioreactors

JoVE 5793

Bioreactors are used to grow organisms in large volumes, thereby enabling the production of mass quantities of the target product. These reactors can be batch reactors, which contain all of the components needed for cell growth, or continuous reactors, which have inlet and outlet ports allowing for the addition of fresh growth media and the removal of cell waste.

This video presents batch and continuous reactors and demonstrates the use of bioreactors to grow bacteria in the laboratory. Finally, this video considers how these reactors are used in the bioengineering field to produce products such as protein therapeutics or even beer.

 Bioengineering

Conformations of Ethane and Propane

JoVE 11708

In an organic molecule, free rotation about the carbon-carbon single bond results in energetically different conformers of the molecule. Due to this rotation, called the internal rotation, ethane has two major conformations — staggered and eclipsed.

Staggered conformation is a low energy and more stable conformation with the C-H bonds on the front carbon placed at 60°dihedral angles relative to the C-H bonds on the back carbon, leading to a reduced torsional strain. In staggered ethane, the bonding molecular orbital of one C-H bond interacts with the antibonding molecular orbital of the other, thereby further stabilizing the conformation. The rotation of farther carbon while keeping the carbon nearer to the observer stationary generates an infinite number of conformers. At  0° dihedral angles, the C-H groups cover one another to form an eclipsed conformation. This conformation has about 12 kJ/mol more torsional strain than the staggered conformation and hence is less stable. Ethane molecules rapidly interconvert between several staggered states while passing through the higher energy eclipsed states. The molecular collisions provide the energy required to cross this torsional barrier.

Similar to ethane, propane also has two major conformers: the stable staggered conformer (low energy) and the unstable eclipsed conformer (higher energy).

 Core: Organic Chemistry

Electrophiles

JoVE 11745

This lesson explains the definition, classification, and characteristic features of an electrophile that are key features of nucleophilic substitution reactions. An analysis of their charge and orbital picture helps understand their reactivity for seeking electrons. Electrophiles can be classified into positive and neutral species. Other classes include free radicals and polar functional groups.

While a positive electrophile, like a proton, reacts due to its vacant, low-energy 1s orbital, the other positive electrophiles, like carbocations, are reactive due to their vacant p orbital.

On the other hand, neutral electrophiles, analogous to Lewis acids, possess empty p orbitals that can accept electrons from the nucleophile to generate stable complexes. An electrophilic center can often be formed in a neutral molecule due to the electron-withdrawing inductive effect in the presence of a more electronegative substituent attached to the molecular chain. This explains the partial positive charge on the carbon atom in a carbonyl group. 

In the context of a chemical reaction, a comprehensive picture of the electron transfer in this process is necessary. A nucleophile deposits its electrons into the lower energy antibonding π orbital of the electrophile. In contrast, the dipole of the σ bond forces the nucleophilic electrons to move into the lower energy antibonding σ orbital, resulting in bond breaking. These phenomena are often explained using examples of carbonyl groups and HCl. Typically, the lowest occupied molecular orbitals (LUMOs) in organic electrophiles are antibonding orbitals with low energy, as they are associated with electronegative atoms. These happen to be either π* orbitals or σ* orbitals.

Some molecules, such as halogens, also make good electrophiles. Here, despite the absence of a dipole, a poor overlap between the atomic orbitals of the two halides weakens the bond making it more prone to a nucleophile attack. This leads to the other classification of strong versus weak electrophiles, covered in later lessons. Usually, molecules with a single or double bond linked to an electronegative atom like O, N, Cl, or Br make good electrophiles.

 Core: Organic Chemistry

Regioselectivity and Stereochemistry of Hydroboration

JoVE 11780

A significant aspect of hydroboration–oxidation is the regio- and stereochemical outcome of the reaction.

Hydroboration proceeds in a concerted fashion with the attack of borane on the π bond, giving a cyclic four-centered transition state. The –BH2 group is bonded to the less substituted carbon and –H to the more substituted carbon. The concerted nature requires the simultaneous addition of –H and –BH2 across the same face of the alkene giving syn stereochemistry.

Figure1

The observed preference in regioselectivity can be explained on the basis of steric and electronic factors.

In the transition state, the larger part of the reagent (–BH2) is bonded to the less substituted carbon, thereby minimizing the steric tension. This results in a less crowded low-energy transition state, which is more stable than Markovnikov's transition state.

Figure2

Further, the addition of borane can result in a partial positive charge on either of the two carbons. However, a partial positive charge on the more substituted carbon is highly favorable, as it gives a more stable transition state. Hence, in order to achieve this, –BH2 must be placed at the less substituted carbon resulting in an anti-Markovnikov orientation.

The second part of the reaction is the oxidation of the product obtained from hydroboration.

Figure3

The migration of the alkyl group in this mechanism occurs with retention of configuration as it transfers with the electron pairs without reconstructing the tetrahedral geometry of the migrating carbon.

Since the reaction is stereospecific, it is essential to recognize the number of chiral centers formed. If one chiral center is formed, both enantiomers are obtained, as syn addition can occur from either face of the alkene with equal probability. However, if two chiral centers are formed, the syn addition dictates which pair of enantiomers is predominantly obtained.

Figure4

 Core: Organic Chemistry

Alkynes to Aldehydes and Ketones: Acid-Catalyzed Hydration

JoVE 11838

Introduction

Analogous to alkenes, alkynes also undergo acid-catalyzed hydration. While the addition of water to an alkene gives an alcohol, hydration of alkynes produces different products such as aldehydes and ketones.       

Figure1

Since the rate of acid-catalyzed hydration of alkynes is much slower than alkenes, a mercuric salt like mercuric sulfate (HgSO4) is usually added to facilitate the reaction. Hydration of terminal alkynes follows Markovnikov's rule; however, for internal alkynes, the addition of water is non-regioselective.

Mechanism

The mechanism begins with a nucleophilic attack by the alkyne π bond on the Hg2+ ion resulting in the formation of a cyclic mercurinium ion intermediate. A second nucleophilic attack by water on the more substituted carbon forms an organomercuric enol that rapidly converts into a stable keto form via keto-enol tautomerism. Protonation of the keto intermediate followed by the loss of an Hg2+ ion yields the enol form of the product. The final step proceeds with the tautomerization of the enol to the desired ketone.

Figure2

Keto-Enol Tautomerism

Unlike alkenes, acid-catalyzed hydration of alkynes is irreversible. This is because the enol intermediate formed during the hydration of alkynes is unstable and rapidly isomerizes to a more stable keto form. The chemical equilibrium that exists between the two forms is referred to as keto-enol tautomerism. Since the C=O bond is considerably stronger than the C=C bond, the equilibrium favors the keto isomer. Keto-enol tautomerism is characterized by the migration of a proton and the change in the location of a double bond.

Acid-catalyzed tautomerization is a two-step process:

Step 1: Addition of proton  across the enol double bond

Figure3

Step 2: Loss of a proton to yield the keto form

Figure4

Example

Acid-catalyzed hydration of 1-propyne initially forms the less stable enol isomer, propen-2-ol, which tautomerizes into a more stable keto product, propan-2-one.

Figure5

Hydration of Terminal And Internal Alkynes

Acid-catalyzed hydration is most useful for terminal and symmetrical internal alkynes because they form only one final product. In contrast, unsymmetrical internal alkynes yield a mixture of products that need to be separated. This lowers the overall yield and makes the process less efficient.

 Core: Organic Chemistry

Structure and Nomenclature of Alcohols and Phenols

JoVE 11920

Overview

Alcohols are one of the most important functional groups in organic chemistry. The name of alcohol comes from the hydrocarbon from which it is derived. Alcohols are organic molecules containing the functional hydroxyl or –OH group directly bonded to carbon. Phenols have an OH group directly attached to a benzene ring. While alcohols are colorless, phenol is a white crystalline compound with a characteristic "hospital smell" odor.

As with other organic compounds, alcohols and phenols are named by formal and standard systems. The most adopted system is the International Union of Pure and Applied Chemistry (IUPAC).

The IUPAC names of the alcohols are derived by adding the suffix 'ol' to the name of the parent alkane. Phenols, on the other hand, are named as hydroxy derivatives of benzene. 'Phenol' is used as the parent name rather than benzene.

Like all organic compounds, alcohols and phenols have several commercial applications. For instance, ethanol is the main ingredient in alcoholic beverages like wine or beer. Phenols are widely used as antiseptics and as disinfectants.

Naming Alcohols

The table below shows the classification and nomenclature of some alcohols and phenols.

Skeletal Structures IUPAC Name Common Name
Figure1 1-butanol n-butyl alcohol
Figure2 2-butanol sec-butyl alcohol
Figure3 2-methyl-2-butanol t-amyl alcohol
Figure4 2-bromo-4-chlorocyclopentanol -
Figure5 2-propen-1-ol allyl alcohol
Figure6 3-cyclohexen-1-ol -
Figure7 1,2-ethanediol ethylene glycol
Figure8 1,2,3-propanetriol glycerol
Figure9 phenol phenol
Figure10 benzene-1,3-diol resorcinol
Figure11 2-methylphenol o-cresol
Figure12 3-methylphenol m-cresol
Figure13 4-methylphenol p-cresol

 Core: Organic Chemistry

Preparation and Reactions of Thiols

JoVE 12112

Thiols are prepared using the hydrosulfide anion as a nucleophile in a nucleophilic substitution reaction with alkyl halides. For instance, bromobutane reacts with sodium hydrosulfide to give butanethiol.

Figure1

This reaction fails because the thiol product can undergo a second nucleophilic substitution reaction in the presence of an excess alkyl halide to generate a sulfide as a by-product.

Figure2

This limitation can be overcome by using thiourea as the nucleophile. The reaction first produces an alkyl isothiourea salt as an intermediate, which forms thiol as a final product upon hydrolysis with an aqueous base.

Figure3

Thiols can readily oxidize to disulfides, sulfinic acid, and sulfonic acid. The oxidation of thiols to disulfides can even occur in the presence of atmospheric air. Thus, the high susceptibility of thiols to undergo air oxidation necessitates the storage of thiols in an inert atmosphere. Oxidation of thiols to disulfides can also be accomplished using reagents like molecular bromine or iodine in the presence of a base. Disulfides, however, can be easily reduced back to thiols by treatment with reducing agents such as HCl in the presence of zinc. Notably, oxidation of thiols to disulfides is a redox reaction. The interconversion between thiols and disulfides is ascribed to the bond strength of the S–S bond, which is approximately half the strength of other covalent bonds.

 Core: Organic Chemistry

Overview of Cell Death

JoVE 12425

Cell death is an essential process where the body gets rid of old or damaged cells. Cell proliferation and death need to be balanced, as an imbalance between the two may lead to cancer or autoimmune diseases.

Cell death was observed in the early 19th century, but there was no experimental evidence to prove it. In 1842, Carl Vogt first discovered cell death in a metamorphic toad; however, it was not termed ‘cell death.’ Scientists discovered different cell death pathways only in the 20th century with advancements in microscopy and histology.

While apoptosis, autophagy, and necrosis are the most commonly known types of death,  some other types are necroptosis, pyroptosis, ferroptosis,entosis, and paraptosis. These types differ in their pathways toward cell death.

Necroptosis is a combination of apoptosis and necrosis. It begins when a ligand binds to death receptors, activating the receptor-interacting protein kinases-1 (RIPK-1) instead of caspase-8,  generally activated during apoptosis. Further, many intermediate proteins are activated, causing cell swelling and lysis. Pyroptosis is a caspase-mediated cell death that involves caspase-1 or caspase-11, required for the maturation of inflammatory cytokines.

In ferroptosis, iron-dependent lipid peroxidation occurs, accumulating reactive oxygen species that eventually induce cell death. Prominent morphological features of ferroptosis include shrinkage of mitochondria and decreased cristae. During entosis, a cell gets internalized into the neighboring cell with the help of Rho GTPase and is then degraded by the lysosomal enzymes.

Paraptosis is activated when the insulin-like growth factor 1 receptor is overexpressed, triggering events leading to cell death. Paraptosis mainly occurs in neurons, and the paraptotic pathway is also being explored in apoptosis-resistant cancer cells. Anoikis is a programmed apoptotic cell death occurring in cells that get detached from the extracellular matrix. These cells follow the extrinsic or intrinsic pathways of apoptosis.

 Core: Cell Biology

Regulation of Hematopoietic Stem Cells

JoVE 12516

All blood and immune cells are produced from the multipotent hematopoietic stem cells (HSCs) by the process of hematopoiesis. However, they all have a limited life span. In addition, many are depleted in immune surveillance or combatting an injury or infection. This makes blood one of the most regenerative tissues. Hematopoiesis helps replenish these blood and immune cells, restoring the body's normal functioning. However, overproduction of blood and immune cells can make them cancerous or activate autoimmune responses. Thus, hematopoiesis needs to be a tightly regulated process.

Regulation of hematopoiesis helps determine whether the HSCs stay quiescent or undergo any other fate such as self-renewal, differentiation, or apoptosis. Several soluble or membrane-bound factors such as cytokines or colony-stimulating factors (CSFs) regulate hematopoiesis.  For example, stem cell factors or steel produced by the bone marrow stroma make contact with c-kit receptors on quiescent HSCs, providing them with survival signals and HSC maintenance. The CSFs are inflammatory cytokines produced by fibroblasts and macrophages in response to allergies and bacterial or parasitic infections. CSFs regulate the rate of HSC proliferation and the number of cell division cycles they must undergo before differentiating into progenitor cells. This ensures the human body possesses an adequate reserve of stem cells to maintain all blood cell types.  CSFs also induce the progenitors to commit to lineage-specific blood and immune cells and perform their functions. They assist with rapid differentiation into specific immune cells and elicit an immediate immune response until the infection is cleared. In the absence of CSFs, HSCs undergo apoptosis.

 Core: Cell Biology

Disassembly of Intermediate Filaments

JoVE 13110

Intermediate filaments (IFs) do not undergo spontaneous disassembly. Enzymes, kinases, and phosphatases add and remove phosphates from specific sites to regulate their disassembly. The IF concentration in the cytoplasm also regulates the disassembly. If the concentration crosses a threshold, it activates the protein kinases in the vicinity, allowing the phosphorylation of IFs.

Keratin proteins, found at the cell periphery near cell junctions, undergo a cycle of assembly and disassembly. In Type III and IV intermediate filaments, the phosphorylation of the N-terminal head domain by the secondary messenger-dependent kinase proteins influences the phosphorylation of the C-terminal tail, that aids in their disassembly.

During mitosis, the transition from prophase to pro-phase results in the nuclear lamins, vimentin, and glial fibrillary acidic proteins undergoing site-specific phosphorylation by Rho kinase, Cdk1, Aurora-B, and PAK1, resulting in disassembly. The phosphorylation of lamins A, B, and C leads to depolymerization of the filaments into lamin dimers, further leading to nuclear membrane disintegration. The lamins remain attached to the disintegrated membrane through their C-terminal prenylation. The removal of phosphates through phosphatase leads to the reassembly of the lamin meshwork during the telophase.

 Core: Cell Biology

Tagging and Fusion Proteins

JoVE 13378

Proteins are involved in several cellular processes and biochemical reactions. Analyzing a specific protein of interest requires it to be isolated from the other proteins in the cell. This is achieved by overexpressing the specific gene in a suitable host to produce large quantities of the target protein. A tag or label is recombined with the gene to produce a fusion protein containing the target protein and the tag. The tags on these fusion proteins can then be used for easy detection and purification processes. Affinity tags, epitope tags, reporter tags, fluorescent tags, and self-splicing intein tags, are just a few of the various protein tags available.

Glutathione S-transferase Tag

Glutathione S-transferase (GST) is a 211 amino acid protein commonly employed to tag recombinant proteins. An expression vector comprising the gene of interest and the DNA sequence for GST is used for expression in a suitable host, such as E.coli. The recombinant protein can be tagged with GST at either its N-terminal or C-terminal. The GST-tag also increases the solubility of the fusion protein compared to the non-tagged native protein. Since GST is an enzyme, it has high binding specificity to its substrate, glutathione. This substrate specificity is used to purify GST-tagged proteins by affinity chromatography using a matrix of glutathione-coated beads.

Tag Cleavage and Self-splicing

While protein tags allow the target protein to be purified, they might hinder further protein analysis. In such cases, the tags are cleaved using proteolytic enzymes. Since these proteases cleave only at specific sites, fusion proteins are designed with such cleavage sites between the target protein and tag. Another method uses self-splicing protein segments, called inteins, that splice the tags from the target protein without additional enzymes. In this method, an intein segment is also recombined into the fusion protein, positioned between the tag and target protein. These inteins self-splice only under certain conditions such as the presence of thiol compounds or specific pH and temperature. Thus, splicing can be specifically induced after purification of the fusion protein to obtain the pure target protein for further analysis.

 Core: Cell Biology

Confocal Fluorescence Microscopy

JoVE 13394

Confocal microscopy is an advanced microscopic technique. The prime advantage of the confocal microscope over other microscopy techniques is its ability to block the out-of-focus light from the illuminated samples using pinholes. It is widely used with fluorescence optics to obtain high-resolution, sharp contrast images. Unlike optical microscopes, confocal microscopes use a focused beam of light laser to scan the entire sample surface at different z-planes. These microscopes are, therefore, very useful for examining thick specimens such as biofilms, which can be examined alive and unfixed.

The confocal microscopes use two pinholes—illumination, and emission pinhole, to modulate the laser beam to obtain clear, crisp images. The laser passes through the illumination pinhole and gets reflected by the dichroic mirror to scan the sample surface. An emission pinhole confocal with the illumination plane; focuses the emitted light reaching the detector. It eliminates light from non-focused z-planes reaching the detector to obtain high contrast two-dimensional images called the optical sections. A computer software program then merges optical sections from different focal planes to reconstruct a three-dimensional image.

Two types of confocal microscope are widely used based on their method of scanning the samples; laser scanning (LSCM) and spinning disc laser (SDLM) microscopy. In LSCM, a point laser scans each focal plane across the sample and collects the emitted fluorescence through a pinhole in detectors. These two-dimensional images, called the optical sections, can be stacked to reconstruct the three-dimensional image. In contrast, SDLM consists of two linked spinning disks with hundreds of pinholes. It allows rapid scanning of the sample surface at different planes and faster image capturing.

Limitations of confocal microscopy

In confocal microscopy, the limited wavelengths of light in the lasers are a disadvantage. The traditional fluorescence microscopes offer a wide range of illumination wavelengths, using mercury or xenon arc lamps as their illumination sources. The high intensity of the laser in early confocal microscopes was damaging for the cells, which has been overcome to a great extent in the multiphoton microscope systems.

 Core: Cell Biology

Whole Body Regeneration

JoVE 13474

Regeneration is the process of restoring injured or lost tissues, organs, or body parts. While simpler organisms generally show greater ability to regenerate their whole body, few complex animals show similarly exceptional regeneration. For example, planarian flatworms have a unique regenerative potential making them a popular study organism among biologists to understand the mechanisms of whole body regeneration. Other organisms, such as hydra, also show extreme regeneration potential; even when cut into several pieces, each piece can turn into an individual organism.

Planarian Stem Cells

Planarians are among the few animals that maintain pluripotent stem cells throughout life. First described in the 1800s, these pluripotent stem cells are now called neoblasts. Neoblasts are essential for the regeneration of the endoderm, mesoderm, and ectoderm, thus, facilitating whole body regeneration. These neoblasts are present throughout the planarian body, including the space surrounding the gut. As much as 30 percent of cells in planaria are neoblasts. Thus, a planarian worm can regenerate a whole body even when cut into small pieces.

Regeneration Polarity in Planarians

Regeneration requires the restoration of complex anatomical structures and specific integration of organs within the body through precise control of the organs’ size, location, and identity. When a planarian flatworm is amputated transversely, two fragments are generated — a fragment containing the head that regenerates a tail; and a tail fragment that regenerates the head. The location along the anterior-posterior axis determines whether the tissue differentiates into the head or tail. The Wnt proteins are a major determinant of this axis, with high Wnt expression in the anterior cells and low Wnt expression in the posterior cells. At the wound site, posterior cells express the Wnt1 gene, which triggers tail regeneration. In contrast, anterior cells express the notum gene, which inhibits the Wnt signaling and facilitates regeneration of the head.

 Core: Cell Biology

Heat Capacities of an Ideal Gas I

JoVE 13684

Heat capacity is the ratio of heat absorbed by the substance corresponding to its temperature change. It is also called thermal capacity and the SI unit of heat capacity is J/K. Whereas, specific heat capacity is defined as the amount of heat necessary to change the temperature of 1 kg of a substance by 1 K and is also called massic heat capacity. Its SI unit is J/kg⋅K.

Molar heat capacity quantifies the ratio of the amount of heat added (or removed) to increase (or decrease) the temperature of 1 mole of a substance by 1 K, either measured under a constant volume or constant pressure. Molar heat capacity is one of the characteristics of a substance and is also called molar specific heat. Its SI unit is J/mol·K. Measuring the molar heat capacity at constant volume is the easiest. The system acquires infinite number of possible heat capacities if neither the pressure nor the volume is constant. In most cases, the molar heat capacity at constant pressure (CP) is higher than the molar heat capacity at constant volume (CV). For example, the air has 40% greater CP than CV.

 Core: Physics

Random Error

JoVE 14508

Random or indeterminate errors originate from various uncontrollable variables, such as variations in environmental conditions, instrument imperfections, or the inherent variability of the phenomena being measured. Usually, these errors cannot be predicted, estimated, or characterized because their direction and magnitude often vary in magnitude and direction even during consecutive measurements. As a result, they are difficult to eliminate. However, the aggregate effect of these errors can be approximately characterized by enumerating the frequencies of observations in a large dataset. Characterizing the collective effect of these errors helps with statistical analyses. Take a large data set of temperature measurements in London, for instance. We can plot the magnitude of temperature vs. the frequency of occurrence, and if the variations (or errors) in the temperature are truly random, we will obtain a normal distribution curve, also known as the Gaussian curve. This curve allows us to apply the mathematical laws of probability to estimate the mean value and the deviation from the mean value, also known as the standard deviation. From there, we can perform tests to eliminate outliers and answer questions about the dataset.

 Core: Analytical Chemistry

Ionic Strength: Effects on Chemical Equilibria

JoVE 14524

The addition of an inert ionic compound increases the solubility of a sparingly soluble salt. For example, adding potassium nitrate to a saturated solution of calcium sulfate significantly enhances the solubility of calcium sulfate. Le Châtelier's principle cannot predict this shift in the equilibrium. Instead, this could be explained in terms of changes in the effective concentration of the ions in solution in the presence of added inert salt.

In this solution, the primary cation—the calcium ion—is surrounded by the primary anion—the sulfate ion—which in turn is surrounded by calcium ions, forming ionic atmospheres around each ion. Additionally, the primary cation and anion are surrounded by the oppositely charged ions of the added inert salt. The ions from the inert compound in the ionic atmosphere causes the net charge on the primary ions—calcium and sulfate, in this case—to decrease, reducing the  frequency of precipitation. This shifts the equilibrium towards the dissociated ion, increasing the solubility of the sparingly soluble salt. This phenomenon is termed the salt effect, electrolyte effect, or diverse ion effect.

The salt effect is highly dependent on the ionic strength of the solution. With an increase in the ionic strength of the solution, more ions diffuse in the ionic atmosphere, causing the net charge on the primary ion to be even lower, facilitating greater dissociation of the salt. Additionally, the charge on the ions constituting the sparingly soluble salt affects the extent of the salt effect. For example, the solubility of doubly charged ions, such as those constituting barium sulfate, is influenced more than the solubility of singly charged ions, such as those constituting silver chloride, by the same concentration of potassium nitrate.

 Core: Analytical Chemistry

Classification of Titrimetric Analysis Based on Reaction Types

JoVE 14540

Titrimetric analysis in solution chemistry involves measuring the volume of solutions and is often called volumetric analysis. The standard solution of known concentration in the burette is called the titrant, whereas the solution of unknown concentration in the flask is called the analyte, or titrand. Titrimetric analyses can be classified into four types based on the reactions between the titrant and analyte.

Titrations between an acid and a base lead to neutralization reactions that form water molecules. This is called an acid-base titration. For example, the titration of a sodium hydroxide standard solution with the analyte hydrochloric acid generates water, leaving behind the acidically neutral sodium and chloride ions. In the second type of titrimetric analysis – complexometric titrations, metal ions such as silver and mercury come together with electron donors such as cyanide and chloride to form complexes. The third type is precipitation titration, in which the titrant reacts with the analyte to form an insoluble product. For instance, the titration between chloride ion and silver nitrate solution will produce the insoluble silver chloride. In the last type of titrimetric analysis, redox reactions are used to determine the amount or concentration of an analyte. In this case, the oxidation states of the titrant and the analyte change as a result of electron transference between the two. The standard solution can be an oxidizing or a reducing agent, and the end point can be detected with indicators or changes in the electrical signals.

 Core: Analytical Chemistry

Precipitation Titration: Overview

JoVE 14579

Precipitation titration involves the reaction of a titrant and an analyte to generate an insoluble precipitate. While precipitation titration uses various precipitating agents, silver nitrate is the most common precipitating reagent; titrations involving Ag+ are called argentometric titrations. Usually, the endpoint in a precipitation titration can be detected by visual indicators.

A precipitation titration curve demonstrates the change in concentration of the titrant or analyte upon adding the titrant. Titration curves are plotted with the volume of the analyte on the x-axis and the p function of the analyte or titrant concentration on the y-axis. A titration curve has three significant regions: before, at, and after the equivalence point. Before the equivalence point, the analyte concentration is in excess. The concentration of the unreacted analyte is calculated from a ratio of the moles of excess analyte in the solution to the total volume of the solution.

At the equivalence point, a stoichiometric amount of titrant has reacted with the analyte to form the precipitate (e.g., AgCl). However, some redissolution of the precipitate gives equal concentrations of Ag+ (titrant) and Cl(analyte). The analyte concentration can be estimated from the square root of the solubility product.

Beyond the equivalence point, the solution contains excess Ag+. Here, the analyte concentration can be estimated using the solubility expression, where the concentration of Ag+ is obtained from the ratio of moles of excess titrant to the total volume.

 Core: Analytical Chemistry

Atomic Nuclei: Types of Nuclear Relaxation

JoVE 14595

Nuclear relaxation restores the equilibrium population imbalance and can occur via spin–lattice or spin–spin mechanisms, which are first-order exponential decay processes.

In spin–lattice or longitudinal relaxation, the excited spins exchange energy with the surrounding lattice as they return to the lower energy level. Among several mechanisms that contribute to spin–lattice relaxation, magnetic dipolar interactions are significant. Here, the excited nucleus transfers energy to a nearby magnetic dipole, usually a tumbling proton.

Spin–lattice relaxation occurs to restore the longitudinal magnetization to its equilibrium value and is characterized by the time constant, T1, which indicates the average lifetime of a nucleus in the excited state. T1 is also called the dipolar or dipole–dipole relaxation time and can range from 0.01 to 100 seconds for liquids. The value of T1 depends on the factors such as the type of nucleus, the location of a nucleus within a molecule, the size of the molecule, and temperature.

Transverse relaxation, also called spin–spin relaxation, occurs when precessing nuclei fall out of phase, resulting in magnetization decay. Transverse relaxation is influenced by static dipolar fields and is usually faster than longitudinal relaxation. The relaxation times observed in typical NMR experiments range from 0.1 to 10 seconds. Additionally, the spin-lattice relaxation time, T1, depends on the applied magnetic field, while T2 is independent of it.

While the relaxation process is essential to prevent saturation and obtain a detectable signal, it also affects the intensity of the NMR signals. Generally, the intensity of the NMR signal is affected by T1 relaxation, whereas shorter T2 results in broadened NMR signals.

 Core: Analytical Chemistry

Inductively Coupled Plasma-Mass Spectrometry (ICP-MS): Interferences

JoVE 14732

Inductively coupled plasma–mass spectrometry (ICP–MS) is a highly selective and sensitive technique for accurate elemental analysis. Though the analysis of ICP–MS mass spectra is comparatively straightforward, it is affected by spectroscopic and non-spectroscopic interferences. Spectroscopic interferences arise when the plasma contains ionic species with an m/z value the same as the analyte ion. Spectroscopic interference can be categorized as isobaric, polyatomic ions, and refractory oxide ion interferences.

Isobaric interference occurs when the inductively-coupled plasma contains isobaric species of analytes. These species are different nuclides with the same mass number and so have the same m/z values as the analyte, assuming the same charge. For example, 40Ar+ and 58Fe+ generally overlap the peaks for 40Ca+ and 58Ni+, respectively. Quadrupole-based mass instruments, having a resolution of less than 1 u, are substantially affected by isobaric interference. However, it can be minimized using higher resolution instruments.

Another spectroscopic interference is observed when various plasma components interact with matrix or atmospheric components to form polyatomic species, which may undergo further fragmentation to form molecular ions having the same m/z value as the analyte ion. 40ArH+, 16O2+, and 40ArO+ are some potential interferents. A blank correction or using a different analyte isotope removes such interferences.

Another critical interference occurs when several analyte and matrix components form oxides and hydroxides whose peaks overlap with the analyte ions. The formation of oxides and hydroxides depends on various experimental factors, like the injector flow rate, the sample orifice size, the sampler skimmer spacing, etc. These interferences are more challenging to eliminate than isobaric and polyatomic ion interferences.

Non-spectroscopic interferences mainly include the matrix effect. It is associated with a high concentration of matrix species—generally, 500 to 1000 mg/mL—causing analyte signal reduction. Specific experimental conditions enhance the analyte signal. Matrix interferences can be eliminated by incorporating an internal standard with mass and ionization potential close to the analyte.

 Core: Analytical Chemistry

The Muscular System

JoVE 14860

The muscular system is essential to the body's overall structure and function, playing a crucial role in movement, stability, and internal processes. It consists of three distinct types of muscle tissue: the skeletal, the smooth, and the cardiac muscles.

  1. Skeletal Muscles: These muscles are under our conscious control and enable movement. They are attached to bones by tendons, and their contraction and relaxation allow us to move. Additionally, they generate heat during physical activity, helping maintain body temperature.
  2. Smooth Muscles: Unlike skeletal muscles, smooth muscles operate without conscious control. They are located in the walls of hollow organs such as the stomach, intestines, blood vessels, and airways. The contraction of these muscles aids in processes like digestion, by propelling food through the gastrointestinal tract.
  3. Cardiac Muscle: Found exclusively in the heart, this type of muscle contracts rhythmically and continuously, pumping blood throughout the body without fatigue.

The Musculoskeletal System

As the name suggests, the musculoskeletal system is a complex network of muscles and bones. It is critical in facilitating movement and providing structural support to the body.

The skeletal component of the system is primarily made up of bones and associated elements like cartilage, tendons, and ligaments. Bones provide our bodies with basic structure and shape and protect vital organs. For instance, the skull shields the brain, while the ribcage safeguards the heart and lungs. The cartilage, tendons, and ligaments further aid mobility and stability by connecting bones and cushioning joints.

The muscular component, on the other hand, consists of various types of muscle tissues. These include skeletal muscles, which are directly attached to bones and are responsible for voluntary movements such as walking or lifting; smooth muscles, which control involuntary actions like digestion and blood flow; and cardiac muscle, found only in the heart, which pumps blood throughout the body. Together, the muscular and skeletal components work harmoniously to allow locomotor functions.

 Core: Anatomy and Physiology

Muscles of the Forearm that Move the Hand and Fingers

JoVE 14878

The muscles of the forearm that move the wrist, hand, and digits are numerous and diverse. They can be classified into two groups based on their location and function — the anterior and posterior compartment muscles.

Anterior Compartment

The anterior compartment muscles originate from the humerus. They primarily function as flexors and are also known as flexor muscles. They typically insert on the carpals, metacarpals, and phalanges. The superficial layer includes the flexor carpi radialis, palmaris longus, and flexor carpi ulnaris. Below these lies the flexor digitorum superficialis, one of the largest superficial muscles in the forearm, which plays a crucial role in flexing the digits. There are two muscles in the deep anterior compartment. The flexor pollicis longus is responsible for flexing the distal phalanx of the thumb. The flexor digitorum profundus, however, ends in four tendons that insert into the distal phalanges of the fingers.

Posterior Compartment

The posterior compartment muscles originate from the humerus and function primarily as extensors. The superficial layer includes the extensor carpi radialis longus, extensor carpi radialis brevis, extensor digitorum, extensor digiti minimi, and extensor carpi ulnaris. The extensor digitorum occupies a significant portion of the posterior surface of the forearm and divides into four tendons that insert into the middle and distal phalanges of the fingers. The extensor carpi ulnaris is located medially in this group. In the deep posterior compartment are four muscles — the abductor pollicis longus, extensor pollicis brevis, extensor pollicis longus, and extensor indicis. These muscles play a crucial role in extending the thumb and fingers.

The Retinaculum

The tendons of these forearm muscles that attach to the wrist or continue into the hand are held close to the bones by strong fasciae. Tendon sheaths surround the tendons, providing protection and reducing friction. At the wrist, the deep fascia forms fibrous bands called retinacula, which hold the tendons in place. The flexor retinaculum is located on the palmar surface of the carpal bones, forming the carpal tunnel where the median nerve and tendons of the flexor digitorum superficialis, flexor digitorum profundus, and flexor pollicis longus pass through. On the dorsal surface of the carpal bones, the extensor retinaculum holds the extensor tendons of the wrist and digits.

 Core: Anatomy and Physiology

Overview of Synapses

JoVE 14895

A synapse is a specialized structure where two neurons connect, allowing them to pass an electrical or chemical signal to another neuron. It is the point of communication between neurons. The term "synapse" is derived from the Greek word "synapsis," which means "conjunction." The entire process of neural communication revolves around the synapse. When activated, a neuron releases chemicals known as neurotransmitters into the synapse. These neurotransmitters cross the synapse and bind to receptors on the receiving neuron, triggering a response. This is how neurons 'talk' to each other.

The significance of synapses in neurobiology is immense. They play a crucial role in the formation of memories and learning. This is due to a phenomenon called synaptic plasticity, which refers to the ability of synapses to strengthen or weaken over time, depending on the amount of activity they experience. This adaptability enables one to learn new things and form new memories.

There are two main types of synapses: electrical and chemical synapses. Electrical synapses allow direct, rapid communication between cells through structures called gap junctions. They are often found in systems that require fast, synchronized activity, for example, in the heart muscle. Chemical synapses, on the other hand, are slower but allow for more complex communication because different neurotransmitters can elicit different responses. They are the most common type of synapse in the human brain.

There are distinct variations between chemical and electrical synapses. In chemical synapses, the transmission of signals relies on the release of neurotransmitter molecules, resulting in a delay of approximately one millisecond between when the axon potential reaches the presynaptic terminal and when the neurotransmitter prompts the opening of postsynaptic ion channels. It is noteworthy that this communication is unidirectional.

Contrarily, electrical synapses facilitate near-instantaneous signaling, which is crucial for synapses involved in reflexes. Some electrical synapses even allow bidirectional signaling. Furthermore, electrical synapses exhibit greater reliability as they are less prone to blockage. These synapses play a crucial role in synchronizing the electrical activity of a group of neurons. For instance, electrical synapses in the thalamus are believed to regulate slow-wave sleep, and their disruption can lead to seizures.

In conclusion, electrical synapses are essential for the proper functioning of many neural circuits. They enable faster information transmission and greater reliability compared to chemical synapses, making them a key component in the processing of signals in the nervous system.

The gap junction is another type of specialized contact between neurons that allows ions and small molecules to pass directly from one neuron to another.

 Core: Anatomy and Physiology

Diencephalon: Hypothalamus and Coordination

JoVE 14911

The hypothalamus is a small yet highly complex and essential brain region that plays a crucial role in regulating various bodily functions. Anatomically, it is located at the base of the brain, just above the brainstem and below the thalamus, forming part of the limbic system.

The hypothalamus interacts with other brain regions, including the pituitary gland, through a direct physical connection called the hypothalamic-pituitary axis. The hypothalamus receives somatic and visceral inputs and regulates various physiological activities within the body. It is divided into the mammillary, tuberal, supraoptic, and preoptic regions. 

The mammillary region comprises the mammillary bodies and the posterior hypothalamic nuclei. This region acts as a relay station in the olfactory pathway and facilitates body temperature regulation.

The tuberal region includes the stalk-like infundibulum, connecting the pituitary gland to the hypothalamus, plus the dorsomedial, ventromedial, and arcuate nuclei. This region contains satiety centers and regulates the pituitary gland's activity.

The supraoptic region contains four important nuclei: paraventricular, supraoptic, anterior hypothalamic, and suprachiasmatic. This region is involved in critical biological functions, such as producing oxytocin and vasopressin and regulating circadian rhythms.

The preoptic region is located at the anterior part of the hypothalamus and contains the medial and lateral preoptic nuclei. This area plays a crucial role in thermoregulation and controlling various bodily autonomic functions.

 Core: Anatomy and Physiology

Cranial and Spinal Meninges

JoVE 14927

The cranial and spinal meninges are complex protective structures surrounding the central nervous system (CNS), consisting of the brain and spinal cord. These meninges consist of the dura mater, the arachnoid mater, and the pia mater. They protect the CNS, provide structural support, and aid in circulating cerebrospinal fluid (CSF).

Cranial Meninges

These meningeal layers cover the cranium. The dura mater is the outermost layer of cranial meninges. It is a thick and durable membrane of dense collagen fibers divided into two layers — the periosteal layer, which attaches to the skull's inner surface, and the meningeal layer. The dura mater also forms dural venous sinuses responsible for draining blood from the brain.

The arachnoid mater is the middle layer. It is a delicate membrane present beneath the dura mater. The subdural space, which contains a small amount of serous fluid, separates the arachnoid mater from the dura mater. The arachnoid mater also has thin strands called trabeculae that connect to the pia mater.

The innermost layer is the pia mater. It is a thin and highly vascularized membrane closely adhering to the surface of the spinal cord and brain. The pia mater contributes to the formation of the blood-brain barrier.

Spinal Meninges

The spinal dura mater is the outermost layer of the spinal meninges. It consists of dense collagen fibers. The dura mater consists of only one layer, known as the meningeal layer. It is separated from the vertebrae by the epidural space, which houses fat and blood vessels.

Beneath the dura mater is the spinal arachnoid mater. The delicate middle layer is separated from the dura mater by the serous fluid-filled subdural space. The subarachnoid space, a CSF-filled space, separates the arachnoid mater from the pia mater.

The spinal pia mater is a thin, vascularized membrane closely adhering to the spinal cord's surface. It extends as a filament called the filum terminale that anchors the spinal cord to the coccyx.

 Core: Anatomy and Physiology

Smooth Muscle Contraction

JoVE 14946

Smooth muscle contraction is a complex process vital for various bodily functions, from maintaining blood vessel tension to facilitating the movement of food through the digestive tract. Unlike striated muscles, smooth muscle contraction begins more slowly and lasts longer.

The onset of contraction is triggered by an increase in calcium ions within the sarcoplasm, similar to the process in striated muscle. However, smooth muscles have a relatively smaller reservoir of the sarcoplasmic reticulum, the primary calcium storage site in striated muscles. Instead, calcium ions infiltrate the smooth muscle sarcoplasm from both the interstitial fluid and the sarcoplasmic reticulum. The absence of transverse tubules, replaced by structures called caveolae in smooth muscle, means that calcium takes longer to reach the central filaments, contributing to the slower onset of contraction typical of smooth muscles.

In the sarcoplasm, the regulatory protein calmodulin binds these calcium ions, activating myosin light chain kinases or MLCK enzymes. These enzymes alter the conformation of myosin heads in the thick filaments through phosphorylation. The ATPase activity of the myosin heads increases, preparing them to form cross-bridges with actin filaments. The myosin ATPases hydrolyze ATP to slide actin filaments over myosin filaments, resulting in a muscle contraction. Unlike skeletal muscles, the calcium removal from smooth muscle sarcoplasm is slow, which keeps the cross-bridges engaged longer, leading to sustained contraction and delaying relaxation.

Smooth muscle also differs from skeletal muscle in that its thick and thin filaments are not arranged into sarcomeres, conferring the ability to stretch considerably while maintaining contractile functionality. When stretched, smooth muscle initially contracts and increases tension but soon undergoes a stress-relaxation response, allowing the muscle to lengthen while reducing tension. This phenomenon, called plasticity, allows the smooth muscle to contract over a range of lengths four times greater than skeletal muscle. Plasticity is crucial in organs like the stomach and bladder, which must accommodate varying volumes without significantly altering internal pressure.

 Core: Anatomy and Physiology

Olfactory Receptors: Location and Structure

JoVE 14965

The process of olfaction, also known as the sense of smell, is a sophisticated chemical response system. The specialized sensory neurons that facilitate this process, known as olfactory receptor neurons, are situated in an upper segment of the nasal cavity, known as the olfactory epithelium. Olfactory sensory neurons are bipolar, with their dendrites extending from the epithelium's apex into the mucus that lines the nasal cavity. Airborne molecules, when inhaled, traverse the olfactory epithelium, dissolving into the mucus. These molecules, called odorants, bind to specific proteins that maintain their mucus solubility and aid in their transport toward the olfactory dendrites. The odorant-protein complexes bind to receptor proteins within the olfactory dendrites' cellular membrane. The receptor proteins are G protein–coupled and generate a graded membrane potential in the olfactory neurons.

The olfactory neuron's axon originates from the epithelial layer's basal surface, traversing an olfactory foramen in the ethmoid bone's cribriform plate and then projects into the brain. The collection of these axons, termed the olfactory tract, interfaces with the olfactory bulb on the frontal lobe's ventral surface. As a result, these axons bifurcate, undertaking diverse paths to various brain locales. Some axons converge on the cerebrum, particularly the primary olfactory cortex in the temporal lobe's inferior and medial regions. Conversely, other target structures are nested within the limbic system and hypothalamus, facilitating the linkage of odors with enduring memory and emotional reactions. An example of this phenomenon is the evocation of emotional recollections by certain scents, such as the aroma of food native to one's place of origin. Notably, olfaction is the singular sensory modality bypassing a synapse in the thalamus before interfacing with the cerebral cortex. This profound interconnection between the olfactory system and the cerebral cortex elucidates why odors can serve as formidable catalysts for memory and emotion.

The respiratory epithelial tissue, including olfactory neurons, may be susceptible to damage from noxious airborne substances. Consequently, olfactory neural cells within the respiratory epithelium undergo periodic regeneration, during which the axons of the newly formed neurons must establish suitable connections within the olfactory bulb. These emerging axons guide their growth pathway by following existing axons in situ within the cranial nerve.

Anosmia: The Impairment of Olfactory Function

The olfactory nerve, pivotal for the perception of smell, may witness a degradation or a complete loss due to severe facial trauma, a scenario frequently observed in vehicular accidents. This particular affliction is referred to as 'anosmia.' The relative movement of the frontal lobe and the ethmoid bone could result in the severing of olfactory tract axons. Individuals engaged in professional combat sports are often susceptible to anosmia due to constant facial and cranial injuries. Moreover, certain medications, notably antibiotics, have the potential to induce anosmia through the extermination of all olfactory neurons simultaneously. The absence of axons within the olfactory nerve implies that the axons from newly generated olfactory neurons lack a pathway to their respective connections in the olfactory bulb. Anosmia can also be transient due to inflammation resulting from respiratory infections or allergies.

Anosmia can diminish the gustatory experience by rendering food tasteless. Individuals with compromised olfactory capacity may necessitate augmented levels of spices and seasonings to detect flavor in their food. There is a potential link between anosmia and mild depressive states, as the diminished pleasure derived from food could potentially instigate a pervasive sense of melancholy.

The olfactory neurons' regenerative capacity diminishes, leading to age-associated anosmia. This can elucidate the heightened use of salt among older adults compared to younger individuals. However, escalated sodium consumption can augment blood volume and arterial pressure, subsequently escalating the probability of cardiovascular diseases among the older demographic.

 Core: Anatomy and Physiology

The Thyroid Gland

JoVE 14981

The thyroid gland is a small, butterfly-shaped gland located in the neck and covers the anterior surface of the trachea. The gland has two lateral lobes connected by a thin tissue mass called the isthmus. Internally, each lobe comprises many small spherical structures known as thyroid follicles, surrounded by a network of blood vessels.

The follicles have a central cavity lined by simple cuboidal to squamous epithelial cells called follicular cells. These cells produce the glycoprotein thyroglobulin and store it in the central colloid. Thyroglobulin acts as a precursor to triiodothyronine or T3 and thyroxine or T4, which are thyroid hormones essential for maintaining the body's metabolic rate and energy production. Additionally, the parafollicular or C cells scattered between the follicles are responsible for producing the hormone calcitonin, which regulates calcium metabolism.

 Core: Anatomy and Physiology

Electric Circuit Elements

JoVE 15066

Circuit elements are the basic building blocks of an electric circuit. Essentially, an electric circuit is the interconnection of these elements. Within electric circuits, one can find two types of elements: passive and active. Active elements have the ability to generate energy, whereas passive elements do not. Passive elements include components like resistors, capacitors, and inductors, while active elements typically encompass generators, batteries, and operational amplifiers.

The most crucial active elements are voltage or current sources, which generally provide power to the connected circuit. There are numerous ways to categorize circuit elements. For instance, distinguishing linear models from nonlinear models is crucial because circuits composed entirely of linear circuit elements are simpler to analyze than those with some nonlinear elements.

An element or circuit is considered linear if the element's excitation and response meet certain conditions. Both superposition and homogeneity properties are satisfied by a linear element. Consider a scenario where the excitation is the current (i) and the response is the voltage (v). If the element is subjected to a current (i1), it delivers a response (v1). Similarly, when the element is exposed to another current (i2), it elicits a response (v2).

Any circuit element that fails to meet either the superposition or the homogeneity principle is classified as nonlinear. Thermistors are examples of such nonlinear elements.

 Core: Electrical Engineering

Cascaded Op Amps

JoVE 15082

Operational amplifiers (op-amps) are versatile electronic components that can be interconnected in a cascade - one after another in a linear sequence. This cascading is possible due to their infinite input resistance and zero output resistance, allowing them to maintain their input-output relationships even when connected in series.

In a cascaded system, each op-amp is referred to as a stage. The output of one stage drives the input of the subsequent stage. As the input signal passes through each stage, it is amplified by the gain of that stage. The resulting overall gain of the cascaded system is equal to the product of the gains of each stage.

This cascading approach offers several advantages. By distributing the total gain among multiple stages, each stage has less gain than if a single op-amp was used. However, since the gain-bandwidth product remains constant for each op-amp, reducing the gain at each stage effectively increases the bandwidth at that stage. This leads to an overall increase in the bandwidth of the cascaded system.

Cascaded op-amps find extensive use in tuned RF amplifiers found within television circuits. These cascaded amplifiers amplify weak signals and improve impedance matching at the input and output, delivering high-quality audio and video signals to the viewers.

However, designing a cascaded op-amp circuit requires careful consideration. The individual gains must be set such that they do not saturate the signal in the various stages of the cascade. Signal saturation can distort the output, leading to poor performance.

 Core: Electrical Engineering

Types of Responses of Series RLC Circuits

JoVE 15099

A second-order differential equation characterizes a source-free series RLC circuit, marking its distinct mathematical representation. The complete solution of this equation is a blend of two unique solutions, each linked to the circuit's roots expressed in terms of the damping factor and resonant frequency.

Equation1

When the damping factor surpasses the resonant frequency, both roots are real and negative, leading to an overdamped response. In this scenario, the circuit's reaction gradually decays over time.

When the damping factor matches the resonant frequency, the second-order differential equation simplifies to a first-order equation with an exponential solution. The natural response follows a pattern of peaking at its time constant and then decaying to zero, signifying critical damping.

Equation2

For situations where the damping factor is less than the resonant frequency, complex roots emerge, characterized by the damped natural frequency. Euler's formula simplifies the complete response to sine and cosine functions, resulting in an underdamped and oscillatory natural response with a time period proportional to the damped natural frequency.

Equation3

These different response behaviors illustrate the significance of source-free RLC circuits in circuit analysis, offering intriguing insights into electrical circuit behavior and applications.

 Core: Electrical Engineering

Op Amp AC Circuits

JoVE 15116

Within an audio system, the filter circuit plays a pivotal role in processing the amplified audio signal from an amplifier. Its primary function is significantly attenuating signal components with lower frequencies, thereby shaping the audio output. This circuit's operations are examined, focusing on the fundamental filter configuration. This configuration involves an operational amplifier arranged in an inverting setup coupled with resistors (R1 and R2) and a capacitor (C1).

Figure1

When faced with a known input signal, the challenge lies in determining the resultant output signal. The first step involves calculating the capacitor's impedance, which is achieved by employing the angular frequency derived from the time-domain expression of the input voltage.

Equation1

As a result, the analysis transitions into the frequency domain, where the input signal is represented in polar form alongside the impedance components Z1 and Z2. Z2 relates explicitly to the parallel combination of capacitor C1 and resistor R2. The core of the analysis rests on applying Kirchhoff's current law and Ohm's law at a specific node in the circuit, thereby formulating a nodal equation for an ideal op-amp. This equation, when rearranged, reveals a critical insight: the ratio of the output to the input voltage is inversely proportional to the ratio of impedances.

Equation2

The known and calculated values are skillfully substituted into this equation to unveil the output voltage in polar form. The outcome, representing the output voltage, can then be transformed into the time domain, providing a comprehensive understanding of the filter circuit's response to the input signal. For the analysis, ideal op-amp properties are often assumed, including the principle that no current enters either of its input terminals and that the voltage across its input terminals remains zero. 

 Core: Electrical Engineering

Buffers: Overview

JoVE 17364

Buffers play a crucial role in stabilizing the pH of a solution by mitigating the effects of small amounts of added acid or base. They consist of a weak acid and its conjugate base or a weak base and its conjugate acid. A solution of acetic acid and sodium acetate is an example of a buffer that consists of a weak acid and its salt: CH3COOH (aq) + CH3COONa (aq). An example of a buffer that consists of a weak base and its salt is a solution of ammonia and ammonium chloride: NH3 (aq) + NH4Cl (aq).

This combination prevents significant pH changes as long as the buffer's capacity is not exceeded. For example, human blood uses a carbonic acid-bicarbonate buffer system to maintain its pH near 7.4. In a buffer, the weak acid component neutralizes added bases by reacting with hydroxide ions, while the conjugate base neutralizes added acids by reacting with hydronium ions.

 Core: Analytical Chemistry

C. elegans Chemotaxis Assay

JoVE 5113

Chemotaxis is a process in which cells or organisms move in response to a chemical stimulus. In nature, chemotaxis is important for organisms to sense and move toward food sources and move away from stimuli that may be toxic or harmful. Chemotaxis is also important at the cellular level. For example, chemotaxis is required for the movement of sperm toward an egg prior to fertilization. In the lab, chemotaxis is frequently examined in the nematode, C. elegans, which is known to migrate towards food sources in soil, but away from toxins such as heavy metals, substances with a low pH, and detergents. This video demonstrates how to perform a chemotaxis assay, which includes preparing the chemotaxis plates and the worms, running the assay, and analyzing the data. Then, we discuss examples of how chemotaxis assays can be used in C. elegans as a tool to understand learning and memory, olfactory adaptation, and neurological disease such as Alzheimer"s disease. Chemotaxis experiments in C. elegans have near-limitless possibilities for learning more about the cellular and genetic mechanisms of many biological processes, and may lead to a greater understanding of human biology, development, and disease.

 Biology I

Mouse Genotyping

JoVE 5160

Even though the human genome was mapped over 10 years ago, scientists are still far from understanding the function of every human gene! One way to evaluate how a gene functions is to disrupt the sequence encoding it and then evaluate the impact of this change (the phenotype) on the animal’s biology. This approach is commonly used in the mouse (Mus musculus), since it shares a high degree of genetic similarity with humans. To track the animals bearing genetic changes over several generations, it is necessary to screen the DNA of each mouse in a process known as genotyping.

This video provides an overview of the theory and practice behind genotyping mice. The discussion begins with the basic principles of mouse genetics, including a review of the terms homozygote, heterozygote, wildtype, mutant, and transgenic. Next, step-by-step instructions are supplied for extracting and purifying genomic DNA from mouse tissue. Examples are provided demonstrating how to interpret genotyping results, as well as how to keep track of mice with the desired genotype. Finally, some representative applications of the genotyping procedure will be presented in order to demonstrate why this common technique is so essential to mouse research.

 Biology II

An Introduction to Cellular and Molecular Neuroscience

JoVE 5213

Cellular and molecular neuroscience is one of the newest and fastest growing subdisciplines in neuroscience. By investigating the influences of genes, signaling molecules, and cellular morphology, researchers in this field uncover crucial insights into normal brain development and function, as well as the root causes of many pathological conditions.

This video introduction to the fascinating world of cellular and molecular neuroscience begins with a timeline of landmark studies, from the discovery of DNA in 1953 to more recent breakthroughs like the cloning of ion channels. Next, key questions in the field are introduced, such as how genes influence neuron activity and how the nervous system is modified by experience. This is followed by brief descriptions of some prominent methods used to analyze genetic material in neurons, manipulate expression of genes, and visualize neurons and their parts. Finally, several applications of molecular and cellular neuroscience are presented to demonstrate how cellular and molecular approaches can be used to profile neuron populations and explore their functions.

 Neuroscience

Fate Mapping

JoVE 5335

Fate mapping is a technique used to understand how embryonic cells divide, differentiate, and migrate during development. In classic fate mapping experiments, cells in different areas of an embryo are labeled with a chemical dye and then tracked to determine which tissues or structures they form. Technological improvements now allow for individual cells to be marked and traced throughout embryonic development and adulthood.

This video reviews the concepts behind fate mapping, and then details a fate mapping protocol in zebrafish using photoactivatable fluorescent proteins. Finally, specific applications and modifications of this unique technique are discussed.

 Developmental Biology

An Introduction to Reward and Addiction

JoVE 5425

Consequences play a major role in controlling our behavior. If the consequence is a reward, then it encourages the associated behavior. Rewards can come in many forms such as a pleasant feeling, money, or food. However, sometimes an individual engages in compulsive behavior despite of negative consequences, and this state is known as addiction. Administration of addictive substances is neurochemically rewarding, which ultimately causes a loss of control in limiting the intake. Scientists aim to better understand the mechanisms behind these concepts and subsequently develop new therapies for treating substance abuse disorders.

JoVE's introduction to reward and addiction explains the neuroanatomical components of the reward pathway. This is followed by some of the important questions asked by behavioral researchers such as how does our brain chemistry change in response to drug use. Prominent methods section reviews some of the tools being employed in the field, like self-administration protocols. Finally, the video discusses example experiments conducted in labs interested in investigating reward and addiction.

 Behavioral Science

Genetic Crosses

JoVE 5541

To dissect genetic processes or create organisms with novel suites of traits, scientists can perform genetic crosses, or the purposeful mating of two organisms. The recombination of parental genetic material in the offspring allows researchers to deduce the functions, interactions, and locations of genes.

This video will examine how genetic crosses were influential in developing Mendel's three laws of inheritance, which form the basis of our understanding of genetics. One genetic crossing technique that was first developed for single-celled organisms such as yeast, known as tetrad analysis, will then be presented in detail, followed by some examples of how this classical tool is used in genetic studies today.

 Genetics

Using Differential Scanning Calorimetry to Measure Changes in Enthalpy

JoVE 5559

Source: Laboratory of Dr. Terry Tritt — Clemson University

Differential Scanning Calorimetry (DSC) is a method of thermodynamic analysis based on heat-flux method, wherein a sample material (enclosed in a pan) and an empty reference pan are subjected to identical temperature conditions. The energy difference that is required to maintain both the pans at the same temperature, owing to the difference in the heat capacities of the sample and the reference pan, is recorded as a function of temperature. This energy released or absorbed is a measure of the enthalpy change (ΔΗ) of the sample with respect to the reference pan.

 General Chemistry

An Introduction to Cell Metabolism

JoVE 5652

In cells, critical molecules are either built by joining together individual units like amino acids or nucleotides, or broken down into smaller components. Respectively, the reactions responsible for this are referred to as anabolic and catabolic. These reactions require or produce energy typically in the form of a “high-energy” molecule called ATP. Together, these processes make up “Cell Metabolism,” and are hallmarks of healthy, living cells.

JoVE’s introduction to cell metabolism briefly reviews the rich history of this field, ranging from early studies on photosynthesis to more recent discoveries pertaining to energy production in all cells. This is followed by a discussion of some key questions asked by scientists studying metabolism, and common methods that they apply to answer these questions. Finally, we’ll explore how current researchers are studying alterations in metabolism that accompany metabolic disorders, or that occur following exposure to environmental stressors.

 Cell Biology

Co-Immunoprecipitation and Pull-Down Assays

JoVE 5695

Co-immunoprecipitation (CoIP) and pull-down assays are closely related methods to identify stable protein-protein interactions. These methods are related to immunoprecipitation, a method for separating a target protein bound to an antibody from unbound proteins. In CoIP, an antibody-bound protein is itself bound to another protein that does not bind with the antibody, this is followed by a separation process that preserves the protein-protein complex. The difference in pull-down assays is that affinity-tagged bait proteins replace antibodies, and affinity chromatography is used to isolate protein-protein complexes.

This video explains CoIP, pull-down assays, and their implementation in the laboratory. A step-by-step protocol for each technique is covered, including the reagents, apparatus, and instruments used to purify and analyze bound proteins. Additionally, the applications section of this video describes a procedure to study how myxovirus proteins inhibit influenza nucleoprotein, an investigation into the role of calcium ions in calmodulin via a pull-down assay, and a modified pull-down assay for characterizing transient protein interactions.

Protein-protein interactions play a significant role in a wide variety of biological functions. The majority of protein-protein interactions and their biological effects have yet to be identified. Co-immunoprecipitation, or CoIP, and pull-down assays are two closely related methods for the identification of stable protein-protein interactions. This video will cover the principles of the two assays, their general laboratory procedures, and applications of these techniques.

CoIP and pull-down assays are variants of immunoprecipitation, a method for selectively isolating a protein species from a complex solution. In an immunoprecipitation experiment, an antibody specific to a target protein is allowed to form an immune complex with that target in the sample. The complex is then captured on a solid support, typically protein A bound to a sepharose bead. Any proteins not captured are removed by centrifugation steps. The protein is then released from the antibody and solid support by boiling in reducing SDS-PAGE sample-loading buffer.

Co-immunoprecipitation is conducted in the same manner, except that intact protein complexes are captured onto the solid support. The antibody binds to the target protein, which, in turn, is bound to another protein that is not targeted by the antibody. As in immunoprecipitation, the protein complex is released from the antibody and solid support by boiling in reducing SDS-PAGE sample loading buffer.

Pull-down assays are similar to co-immunoprecipitation, differing only in the use of a "bait" protein, as opposed to an antibody. Through molecular biology techniques, this bait protein is engineered with an affinity tag, such as a series of histidine residues. These affinity tags are described in "Chromatography-based Biomolecule Separations." The protein is then captured on an immobilized affinity ligand specific for the tag. The captured protein is then incubated with a sample that contains proteins that form complexes with the "bait". The protein complex is released from the affinity support by washing with a solution containing a competitive analyte specific for the tag on the "bait" protein. Pull-down assays are useful for confirming protein interactions predicted by co-immunoprecipitation, and for discovering unknown protein interactions.

Now that the principles of co-immunoprecipitation and pull-down assays have been discussed, let's look at their laboratory procedures.

First let's discuss co-immunoprecipitation. To a series of microfuge tubes, the following is added: PBS buffer, and a 50% solution of protein A-sepharose, a resin-protein complex that binds to the antibody. The microfuge tubes are rotated to ensure proper distribution, and then the resin is washed with additional PBS buffer. Cell lysate, containing the desired protein, and 2 μg of antibody are added to the microfuge tubes, and the mixture is rotated for 1 h at 4 °C. The beads are pelleted by centrifugation, the supernatant is discarded, and the beads re-washed three times with buffer to remove non-bound protein. The beads containing the antibody-protein complex are in reducing SDS-PAGE sample-loading buffer, for removal of the complex from the antibody and for analysis by SDS-PAGE and immunoblotting.

Now we will discuss the procedure for pull-down assays: The "bait" protein is expressed in a plasmid with the appropriate affinity tag. After reaching log-phase growth, the cells are lysed, and then centrifuged. Suspended streptavidin-sepharose beads, which capture biotin-tagged "bait" protein, are pipetted into a microfuge tube. The beads are then centrifuged and the supernatant carefully removed by aspiration. The beads are then washed with buffer, centrifuged, and the supernatant removed.

Cells containing the putative "prey" protein, which has an affinity for the "bait" protein, are harvested by centrifugation. The supernatant is then added to the microfuge tube containing the resin, and incubated at 4 °C for 3 h on a shaker. The resin is then centrifuged, the supernatant removed, and the resin washed to remove non-bound protein. Elution buffer is added to the resin, and the mixture is incubated at room temperature for 30 min on a shaker. The resin is then centrifuged, and the supernatant containing the desired complex is analyzed by immunoblotting.

Now that we've reviewed the procedures, let's look at some of the useful applications of co-immunoprecipitation and pull-down assays.

Co-immunoprecipitation can be useful in better understanding of enzymes' mechanism of action. Myxovirus-, or Mx-, resistance proteins inhibit a wide range of viruses, including influenza A, for which the mechanism is poorly understood. Co-immunoprecipitation was used to study the interaction between mouse Mx1 protein and influenza nucleoprotein.

Pull-down assays have proven useful in studying the effects of second messengers, which are proteins that communicate a signal from the cellular environment. They are a component of a signaling pathway where multiple proteins interact in response to environmental cues. Calcium ions act as secondary messengers by binding to calmodulin, which, in turn, binds to a wide variety of proteins, mediating many types of biological responses. Without the calcium, the protein can't bind to the calmodulin. A pull-down assay was conducted to test the ability of proteins to bind to calmodulin in the presence or absence of calcium ions.

Co-immunoprecipitation and pull-down assays are generally used for analyzing stable or strong protein interactions, but not transient ones. A recent development in pull-down assays, the HaloTag, has simplified the study of transient protein interactions; HaloTag is a genetically-encoded protein fusion tag, fused to the protein of interest, capable of chemically reacting with a haloalkane solid support. If functional analysis is desired, the protein complex, minus the tag, in its entirety could then be isolated by incubating with tobacco etch virus protease.

You've just watched JoVE's video on co-immunoprecipitation and pull-down assays. This video described the principles of the two methods, the general laboratory procedures, and some of their applications.

Thanks for watching!

 Biochemistry

Overview of Biosensing

JoVE 5794

Biosensors are devices that use a wide range of biological processes and physical properties in order to detect either a biological molecule, such as a protein or cell, or a non-biological molecule, such as a chemical component or contaminant. This interdisciplinary field utilizes electrical, optical, electrochemical, or even mechanical properties to detect the presence of the target molecule.

This video introduces the field of biosensing, and reviews common types of biosensor technologies. This video also discusses key challenges in the field, and provides insight into how biosensors are used in the field.

 Bioengineering

Conformations of Butane

JoVE 11709

Unlike ethane and propane that have only two major conformations, butane has more than two conformers. The staggered form of butane in which the bulky methyl groups on the two carbons are placed on opposite sides, that is, at a dihedral angle of 180°, is the lowest energy, most stable form — called the anti conformer. This conformation is stabilized due to the absence of steric repulsion between the largely spaced out methyl groups. The other two staggered conformations are degenerate and have 3.8 kJ/mol more energy than the anti conformation, emerging from the steric repulsion between the methyl groups that are now positioned at 60° dihedral angles. These unfavorable steric interactions are known as gauche interactions, and the conformers, as such, are called gauche conformers.

The totally eclipsed butane, due to the two eclipsing methyl groups at 0° dihedral angle, has the highest energy (19 kJ/mol) and is the most unstable form. The other two eclipsed conformations — with two CH3-H eclipses and one H-H eclipse — are degenerate, each with an energy cost of 16 kJ/mol.

 Core: Organic Chemistry

SN1 Reaction: Mechanism

JoVE 11753

Kinetic studies of ionization of a tertiary halide in a protic solvent suggest that only the substrate participates in the rate-determining step (slow step). The nucleophile is involved only after the slowest step. The SN1 reaction takes place in a multiple-step mechanism. 

Firstly, the haloalkane ionizes to generate a carbocation intermediate and a halide ion. This heterolytic cleavage is highly endothermic with large activation energy. The ionization of the substrate, facilitated by a polar protic solvent, is the slowest of all steps, making it the rate-determining step of an SN1 reaction. The ions formed are stabilized through solvation. In the second step, the reactive carbocation intermediate behaves as a strong electrophile and is attacked by the nucleophilic solvent molecule that quickly donates an electron pair to generate an oxonium ion. This process is exothermic. In the third step, the solvent abstracts a proton from the oxonium ion to yield the final nucleophilic substituted product. 

Thus, the SN1 reaction consists of two core steps for substitution and an additional step of proton loss. The mechanism further suggests that several factors such as the stability of the carbocation, the nature of the leaving group, and the nature of the solvent used, favor the SN1 mechanism. 

 Core: Organic Chemistry

Oxidation of Alkenes: Syn Dihydroxylation with Osmium Tetraoxide

JoVE 11781

Alkenes are converted to 1,2-diols or glycols through a process called dihydroxylation. It involves the addition of two hydroxyl groups across the double bond with two different stereochemical approaches, namely anti and syn. Dihydroxylation using osmium tetroxide progresses with syn stereochemistry.

Figure1

Syn Dihydroxylation Mechanism

The reaction comprises a two-step mechanism. It begins with the addition of osmium tetroxide across the alkene double bond in a concerted manner forming a five-membered cyclic osmate ester as an intermediate, which can be isolated and characterized. Osmium tetroxide is electrophilic in nature, serving as a strong oxidizing agent. It accepts an electron pair from the alkene π bond undergoing a reduction from +VIII to +VI.

In the next step,  the cyclic osmate ester reacts with a reducing agent like sodium bisulfite that cleaves the Os–O bond producing a cis-glycol with retention of the syn stereochemistry of the two newly formed C–O bonds.

Figure2

A major drawback of the method is the use of toxic and expensive osmium tetroxide. To overcome this, osmium tetroxide is often used as a catalyst along with the co-oxidants like N-methylmorpholine N-oxide (NMO) or tert-butyl hydroperoxide (TBHP). The co-oxidants reoxidize the osmium +VI species to +VIII, thereby regenerating osmium tetroxide for further oxidation of the remaining alkenes.

Stereochemical Outcome

As oxidation of alkenes using osmium tetroxide is a stereospecific syn addition process, the two oxygens of osmium tetroxide are simultaneously added to the same face of the alkene π bond. Based on this, dihydroxylation of (E)-hex-3-ene produces a pair of enantiomers, while (Z)-hex-3-ene gives a meso compound.

Figure3

Sharpless Asymmetric Dihydroxylation

Interestingly, Karl Barry Sharpless developed an enantioselective method for syn dihydroxylation of alkenes, for which he was awarded the Nobel prize. This method is known as Sharpless asymmetric dihydroxylation that is carried out using osmium tetroxide, a stoichiometric amount of the co-oxidant, and a chiral amine ligand. 

Figure4

 Core: Organic Chemistry

Alkynes to Aldehydes and Ketones: Hydroboration-Oxidation

JoVE 11839

Introduction

One of the convenient methods for the preparation of aldehydes and ketones is via hydration of alkynes. Hydroboration-oxidation of alkynes is an indirect hydration reaction in which an alkyne is treated with borane followed by oxidation with alkaline peroxide to form an enol that rapidly converts into an aldehyde or a ketone. Terminal alkynes form aldehydes, whereas internal alkynes give ketones as the final product.

Figure1

Mechanism

The hydroboration-oxidation reaction is a two-step process. It begins with the hydroboration step, which involves a concerted syn addition of BH3 across the carbon–carbon triple bond to form an alkenylborane. The concerted nature of the reaction also accounts for the anti-Markovnikov regiochemistry, where the BH2 group adds to the less substituted carbon and H to the more substituted carbon of the triple bond.

Figure2

 Three successive hydroboration reactions convert an alkene into a trialkenylborane intermediate. The second part of the sequence is oxidation, where the trialkenylborane is treated with alkaline hydrogen peroxide to form an enol. The enol eventually converts into a stable carbonyl product via keto-enol tautomerism.

Figure3

Hydroboration of Alkynes with Disubstituted Boranes

Unlike alkenes, hydroboration of alkynes does not stop at the first addition of BH3. This is because alkynes have two π bonds, each capable of reacting with BH3. The first addition forms an organoborane, which is an alkene derivative that can react further with another equivalent of BH3.

Terminal alkynes being less hindered than internal alkynes are more susceptible to a second BH3 addition. With internal alkynes, the addition of BH3 stops after the first stage and proceeds in a direction to give the trialkenylborane.

Figure4

 Nevertheless, hydroboration of terminal alkynes can be stopped at the first step by using bulky disubstituted boranes (R2BH) such as disiamylborane and 9-BBN instead of BH3.

Figure5

The first addition of the bulky reagent forms a sterically hindered alkenylborane that resists any further additions and helps in the efficient conversion of alkynes to stable carbonyl compounds.

Figure6

 Core: Organic Chemistry

Physical Properties of Alcohols and Phenols

JoVE 11921

Alcohols are organic compounds in which a hydroxy group is attached to a saturated carbon. Phenols are a class of alcohols containing a hydroxy group attached to an aromatic ring. The physical properties of the alcohols and phenols are influenced by hydrogen bonding due to the oxygen–hydrogen dipole in the hydroxy functional group and dispersion forces between alkyl or aryl regions of alcohol and phenol molecules.

Alcohols possess a higher boiling point than aliphatic hydrocarbons of similar molecular weights due to intermolecular hydrogen bonding. As in hydrocarbons, the dispersion forces are the reason for the higher boiling point upon increasing carbon chain length.

Hydrogen bonding between the hydroxy group and water facilitates the solubility of alcohols in water. However, water solubility also depends on the length of the alkyl or nonpolar region of the molecule. Alcohols with an alkyl region of up to three carbon atoms are miscible with water. As the chain length increases, the increased surface area of the nonpolar region hinders solvation by the water molecules.

The solubility of branched alcohols is higher than that of linear alcohols of similar molecular weight. Branching reduces the surface area for intermolecular interactions between nonpolar regions; hence, the hydrophobic nonpolar region is smaller. Because of the weaker intermolecular interactions, the boiling points of branched alcohols are lower than the corresponding linear alcohols.

Multiple sites for hydrogen bonding in one molecule increase the boiling point; therefore, diols and amino alcohols possess higher boiling points and better water solubility than alcohols.

Compared to linear alcohols, cyclic alcohols can only exist in a limited number of conformations due to steric restrictions. The increased intermolecular interactions that arise from the close packing of cyclic alcohol in the liquid phase results in a higher boiling point as compared to that of a linear alcohol.

Intermolecular hydrogen bonds also play a role in defining phenols' high boiling point and solubility in water. The boiling point of phenol is higher than that of the corresponding aliphatic alcohol due to the close packing of phenol molecules, facilitated by π–π stacking interactions between the large, planar aromatic rings. Close-packed aromatic rings increase the surface area of the nonpolar region in the liquid phase and limit the solubility of phenol (9.3 g in 100 g H2O). However, this solubility is higher than that of alcohols with a similar molecular weight due to the increased polarity of the oxygen–hydrogen bond dipole induced by adjacent electron-withdrawing aromatic rings.

Structure Name Molecular weight (g/mol) Boiling point (oC) Solubility

(g/100 g H2O)

Figure1 1-Butanol 74 118 9.1
Figure2 Isobutanol 74 108 10
Figure3 tert-Butanol 74 83 miscible (∞)
Figure4 Pentane 72 36 insoluble
Figure5 Propane-1,2-diol 76 188 miscible (∞)
Figure6 1-Hexanol 102 156 0.6
Figure7 Cyclohexanol 100 162 3.6
Figure8 Phenol 94 182 9.3
Figure9 Toluene 92 110 insoluble

Alcohols are widely used as antiseptics due to their antibacterial properties. Isopropanol or ethanol is the major component of hand sanitizer. An ideal antibacterial agent should have a significant nonpolar region or alkyl region that can penetrate the cell membranes of microorganisms and destroy them. At the same time, it should have high solubility in the transport medium, which is water. In smaller alcohols, the optimal balance between these two conditions is fulfilled.

 Core: Organic Chemistry

Preparation and Reactions of Sulfides

JoVE 12113

Sulfides are the sulfur analog of ethers, just as thiols are the sulfur analog of alcohol. Like ethers, sulfides also consist of two hydrocarbon groups bonded to the central sulfur atom. Depending upon the type of groups present, sulfides can be symmetrical or asymmetrical. Symmetrical sulfides can be prepared via an SN2 reaction between 2 equivalents of an alkyl halide and one equivalent of sodium sulfide.

Figure1

Asymmetrical sulfides can be synthesized by treating thiols with an alkyl halide and a base. The reaction follows an SN2 pathway and proceeds via the thiolate ion intermediate. This reaction is the sulfur analog of Williamson ether synthesis and prefers methyl, primary, and secondary alkyl halides but not tertiary alkyl halides. Sulfides can readily oxidize to sulfoxide and sulfones.

Figure2

Treatment of sulfide with one equivalent of hydrogen peroxide at room temperature yields sulfoxide, which upon further oxidation with a peroxy acid yields a sulfone. However, 2 equivalents of hydrogen peroxide oxidize sulfide directly to sulfone.

Figure3

Dimethyl sulfoxide (DMSO) and tetramethylene sulfone are common examples of sulfoxides and sulfones, respectively. They both are excellent dipolar aprotic solvents.

 Core: Organic Chemistry

Apoptosis

JoVE 12426

Apoptosis is a combination of two Greek words, 'apo' and 'ptosis,' meaning separation and falling off, respectively. Hippocrates used this word to describe gangrene, which was caused due to bandaging of fractured bones. Apoptosis was distinguished from necrosis in 1970 when John Kerr reported observations of morphological changes occurring during apoptosis. During one experiment, he observed that the disruption of blood supply to the liver tissue resulted in a size reduction of the tissue. After examining the tissue under an electron microscope, he observed that the hepatocytes had shrunk and the chromatin had condensed. In 1980, Wyllie established the relationship between DNA degradation and apoptosis.

Apoptosis allows a cell to die in a controlled manner that prevents the release of potentially damaging molecules from inside the cell to the extracellular space. Many internal checkpoints monitor a cell's health. If abnormalities are observed, a cell can spontaneously initiate apoptosis. However, the cell's standard checks and balances fail in some cases, such as during viral infection or uncontrolled cell division due to cancer. External signaling can also initiate apoptosis. For example, most normal animal cells have receptors that interact with the extracellular matrix, a network of glycoproteins that provides structural support for cells in an organism. The binding of cellular receptors to the extracellular matrix initiates a signaling cascade within the cell. However, if the cell moves away from the extracellular matrix, the signaling ceases, and the cell undergoes apoptosis. This system prevents cells from traveling through the body and proliferating out of control, as happens with tumor cells that metastasize.

Regulated apoptosis is essential to maintain normal physiology and tissue homeostasis. Dysregulated apoptosis in cells leads to various diseases, such as cancer, Alzheimer's, and cardiovascular diseases. Dysregulated apoptosis can also cause increased apoptosis, as observed during various autoimmune diseases like Hashimoto's thyroiditis, systemic lupus erythematosus, and rheumatoid arthritis.

Some portion of this text is adapted from Openstax, Biology 2e, Section: 9.3

 Core: Cell Biology

Mesenchymal Stem Cells

JoVE 12518

Mesenchymal stem cells (MSCs) are adult stem cells that can differentiate into most connective tissue cell types, except for hematopoietic cells, depending upon the source of MSCs. For example, bone-marrow-derived MSCs (BM-MSCs) can differentiate into osteocytes, hepatocytes, and pancreatic and neuronal cells. MSCs can be isolated from various sources such as bone marrow, placenta, adipose tissue, teeth, and Wharton’s jelly, a gelatinous substance in the umbilical cord. The ease of their access and in vitro expansion and stability have made MSCs the choice of cells in regenerative medicine.

MSCs also show anti-inflammatory and immunomodulatory properties. Immunomodulation is the ability to activate or suppress the immune system to respond to an infection. Immunomodulation is facilitated by extracellular vesicles (EVs) released from MSCs. EVs are reported to efficiently carry molecules such as nucleic acids, proteins, and other paracrine signaling molecules to recipient cells. The EVs can also cross the blood-brain barrier and effectively trigger immune responses. Hence, MSCs and EVs are being investigated for their use as therapeutic agents to treat diseases such as Rheumatoid arthritis and other autoimmune disorders.

While MSCs are explored for their therapeutic potential in repair and regeneration, MSCs also show pro-tumorigenic properties. MSCs secrete growth factors that induce angiogenesis and facilitate tumor growth and invasion, especially when close to other tumor cells. Despite this challenge, MSCs are used as therapeutic agents for cancer treatment. MSCs can inhibit pathways related to the tumor cell cycle or release cytotoxic factors to induce programmed cell death in cancer cells. Thus, MSCs are a ‘double-edged sword’ in cancer treatment.

 Core: Cell Biology

Flow Cytometry

JoVE 13379

The development of flow cytometry techniques began in 1934 with initial attempts by Andrew Moldavan, a bacteriologist who counted the cells in a flowing capillary system. Moldavan pumped cells through a capillary tube focused under a microscope for visualization. The invention of photometry allowed the measurement of differentially-stained cells, and Louis Kamentsky developed the first multiparameter flow cytometer in 1965 to identify and count the cancer cells in cervical tissue specimens.

In 1967, Kamentsky and Melamed achieved the diversion of cells away from the flow, allowing cell sorting. But it was not until 1973 that a refined electrostatic sorter, known as a fluorescence assisted cell sorter or FACS, was developed. Typically, the cell suspension passes through a laser beam, and the light scatter, or fluorescence, is detected. The changes in scatter patterns or fluorescence determine how the cells should be sorted.

Applications of Flow Cytometry

Apart from cell counting and sorting, flow cytometry is also used to estimate the DNA and RNA content of cells. When a cell suspension is incubated with fluorochromes, viable cells take up the stains to different degrees based on cell type and growth phase. These differences are detected by the FACS software, which then estimates the relative proportions of the cell types in different growth phases. Since flow cytometry estimates cell parameters such as shape, size, and granularity, it can also be used to differentiate necrotic and apoptotic cells. For example, apoptotic cells are typically more granular and therefore tend to have a higher side scatter of light than the less granular necrotic cells.

Limitations of Flow Cytometry

An obvious limitation of flow cytometry is that it cannot analyze cells that are not suspended in culture, such as those growing on an adherent medium or organized in a tissue. All samples must be dissociated into cell suspensions to generate single-cell droplets; therefore, cell-cell interactions cannot be analyzed. Another constraint stems from the limit on the number of simultaneous detectors. Despite many fluorophores being available, a flow cytometer usually does not have more than 12 detectors to efficiently and consistently detect the different cell subpopulations. While flow cytometry is a powerful technique, it can also generate an overwhelming quantity of data that requires a human expert to analyze and process. And finally, given the high sensitivity of this technique, standard protocols for sample preparation, data recording, and analysis are essentia. However, a lack of standardization of these aspects makes it  difficult to compare results from different studies.

 Core: Cell Biology

Protein Dynamics in Living Cells

JoVE 13395

Different fluorescence-based techniques are used to study the protein dynamics in living cells. These techniques include FRAP, FRET, and PET.

Fluorescent recovery after photobleaching (FRAP) is a fluorescent-protein-based detection technique used to quantify protein movement rates within the cell. This method exposes a small portion of the cell to an intense laser beam. The laser beam causes permanent photobleaching of the fluorophore-tagged proteins in the exposed region. As the bleached proteins diffuse out over time and fluorescently-labeled proteins from other parts of the cell move in, the region becomes fluorescent again. The rate of movement is quantified by plotting the relative fluorescence intensity versus the time taken, which allows for determining the rate with which protein moves within the cell.

Förster resonance energy transfer (FRET) is a molecular technique used to determine the distance between two proteins. Thus, FRET can be used as a molecular ruler. The technique uses two fluorophore-tagged proteins, one acting as the donor and the other as the acceptor. Upon excitation at a specified wavelength, the donor fluorophore emits fluorescence energy that is absorbed by the acceptor fluorophore. The donor fluorophore returns to the ground state upon energy transfer, while the acceptor fluorophore emits fluorescence that is visualized with a fluorescence microscope. The FRET technique depends on three factors, the distance between interacting proteins, the extent of spectral overlap between the donor and acceptor fluorophores, and the orientation of the donor and acceptor fluorophore during energy transfer. The energy transfer between the correctly oriented donor and acceptor fluorophores can only occur when the distance between the two interacting proteins is 10 nm or less.

Photoinduced Electron Transfer or PET determines the sub-atomic distance between proteins in a cell. In PET, the fluorophore absorbs the light and emits a fluorescence signal through an excited electron. The excited electron is transferred to the receptor. During the energy transfer, a redox reaction occurs, generating charge separation between the donor and acceptor proteins.

 Core: Cell Biology

Stem Cell Therapy for Tissue Regeneration

JoVE 13475

Stem cell therapy is a method used in regenerative medicine to repair and restore function to damaged tissues and organs. Stem cells have the potential to proliferate and differentiate into various tissue types, making them ideal candidates for tissue regeneration. For example, hematopoietic stem cell transplants are commonly used in blood cancer treatment to replenish damaged bone marrow and restore healthy blood cells.

Types of Stem Cells used in Stem Cell Therapy

The two main cell types that have the potential for use in stem cell therapy are adult stem cells and embryonic stem cells (ESCs). However, ESCs, derived from the blastocyst, require the sacrifice of embryos, which raises several moral and ethical issues. Hence, adult stem cells, such as hematopoietic (HSCs) and mesenchymal stem cells (MSCs), are more commonly used for regenerative therapy despite their limited differentiation potential.

HSCs and MSCs are both primarily found in the bone marrow of adults and the umbilical cord blood. While HSCs can differentiate and produce different blood cells, including RBCs and WBCs, MSCs can differentiate into tissues of the heart, liver, muscle, bone, and also some nerves.

Strategies in Stem Cell Therapy

Stem cell therapy is a complex process that requires optimization at each stage. Optimization depends on many factors, such as the source of stem cells, the target tissue, and the recipient's health condition. For example, intravenously administered HSCs can establish and proliferate in the bone marrow. In contrast, direct administration of stem cells at the target site is more efficient for the local repair of tissues in organs such as the heart, kidney, liver, and muscle. In such tissues, the timing of stem cell introduction also plays a crucial role in determining repair efficiency. When stem cells are administered too early, such as just after the injury occurs, the inflammatory process impairs the proliferation of these stem cells. However, if stem cells are introduced during the repair phase, they are more readily incorporated into the tissue and can proliferate to repair the damage.

 Core: Cell Biology

Force On A Current Loop In A Magnetic Field

JoVE 13770

Magnetic forces on wires carrying current are most frequently applied in motors. A DC motor is a device that converts electrical energy into mechanical work. In motors, wire loops are enclosed in a magnetic field. When current flows through the loops, the magnetic field applies torque, which causes the shaft to rotate. The direction of the current is reversed once the loop's surface area is lined up with the magnetic field, causing a constant torque on the loop. During the process, commutators and brushes are used to reverse the current. The commutator is set to reverse the current flow at certain points to keep continual motion in the motor. A basic commutator has three contact areas to avoid and dead spots where the loop would have zero instantaneous torque at that point. The brushes press against the commutator, creating electrical contact between parts of the commutator during the spinning motion.

Consider a rectangular-shaped current-carrying loop of wire with sides of lengths a and b, such as a loop in a motor, placed in a uniform magnetic field, as shown in figure 1.

Equation1

The loop of wire experiences forces that can be calculated by applying the equation for the magnetic force on a straight current-carrying wire to each of the sides. The force on side 1 and side 2 can be calculated as follows:

Equation1

Equation2

The magnitude of the current on side 3 and side 4 is the same but flows in the opposite direction to that on side 1 and side 2, respectively. Thus, the force on side 3 and side 4 is equal to the force on side 1 and side 2, respectively, with a negative sign. Finally, the total force in the current-carrying loop of wire is equal to the sum of the individual forces acting on each side of the loop and can be calculated as follows:

Equation3

 Core: Physics

Standard Deviation of Calculated Results

JoVE 14509

Standard deviation measures the spread of data around the mean value. Many large data sets follow a Gaussian distribution, also known as a normal distribution. This distribution is bell-shaped curved, with the most frequently observed value (mean or central value) in the middle. The farther away from the central value, the greater the deviation from the central value, and the lower the frequency.

A broad Gaussian distribution curve has a wider standard deviation, representing a data set with lower precision. For data with high precision, the magnitude of the standard deviation is small, and the width of the distribution curve is narrow.

The mean and standard deviation of the entire data or population are known as the population mean and population standard deviation, respectively. The standard deviation for a population subset is the sample standard deviation, and the estimated mean of the subset is known as the sample mean. A useful way to express the standard deviation is the 'relative standard deviation', which expresses the standard deviation as a fraction of the mean. This quantity is also known as the coefficient of variation and is often expressed as a percentage. 

 Core: Analytical Chemistry

Thermodynamics: Chemical Potential and Activity

JoVE 14525

The effective concentration of a species in a solution can be expressed precisely in terms of its activity. Activity considers the effect of electrolytes present in the vicinity of the species of interest and depends on the ionic strength of the solution. The activity of a species is expressed as the product of molar concentration and the activity coefficient of the species.

The thermodynamic equilibrium constant is more accurately defined in terms of activity rather than concentration. Activity is formally defined in terms of the chemical potential, also known as the partial molar Gibbs energy. The chemical potential of a system is the Gibbs energy change of the system per mole of a species in question, given that the temperature, pressure, and the number of moles of other species are held constant. For a chemical reaction to take place spontaneously, the chemical potential of the products must be less than that of the reactants. If the number of moles of other species varies in the solution—i.e., a change in the composition of the solution—the chemical potential of the solution changes, and so does the activity.

 Core: Analytical Chemistry

Titration of a Weak Base with a Strong Acid

JoVE 14541

The titration curve of a weak base like ammonia with a strong acid like hydrochloric acid is the mirror image of the titration curve of a weak acid with a strong base.

Using the ICE table and substituting the Kb value, we calculate the initial pH of 50 mL of 0.1 M ammonia to be 11.11. Addition of 25 mL of 0.1 M hydrochloric acid to this solution of ammonia results in a buffer with an equal concentration of ammonia and ammonium ions. The pH of this buffer can be calculated by substituting these values into the Henderson-Hasselbalch equation, which shows that at this point, also known as the half-equivalence point, the pH is equal to pKa. Continuing the titration to the equivalence point would mean that a total of 50 mL of 0.1 M hydrochloric acid would have been added, converting all the ammonia molecules to ammonium ions. At this point, the ammonium ions undergo hydrolysis to generate hydronium ions, making the solution acidic, and an ICE table calculation shows that the pH is 5.28. After the equivalence point, as the hydronium ions predominate the solution, the pH further drops to 1.6. The titration endpoint is detected using the indicator methyl red, which exhibits a color change in the desired pH range.

 Core: Analytical Chemistry

Precipitation Titration Curve: Analysis

JoVE 14580

The precipitation titration curve demonstrates the change in concentration of one reactant with the volume of titrant added. During the titration of chloride ions with silver nitrate, the precipitation titration curve is divided into three regions: before, at, and after the equivalence point. Before the equivalence point, low redissolution of the sparingly soluble silver chloride precipitate gives a low silver ion concentration. However, in the second region, representing the equivalence point, the silver ion concentration sharply increases as chloride ions are completely consumed. In the last region, beyond the equivalence point, the silver ion concentration is high, mainly due to excess titrant.

The shape of the curve is influenced by the reactant concentration and the solubility product of the precipitate, as revealed by comparing the titration curves of the three halides. Because the solubility product of AgI is smaller than those of AgBr and AgCl, the least soluble AgI precipitates first, indicating a larger break at the equivalence point.

So, if a mixture of KI and KCl is titrated against AgNO3, AgI precipitates first. After the iodide ions are entirely consumed, the silver ion concentration increases, and AgCl precipitates. After completely consuming the chloride ions, the silver ion concentration increases again, and two equivalence points are observed.

 Core: Analytical Chemistry

Chemical Shift: Internal References and Solvent Effects

JoVE 14596

In an NMR sample, precise measurement of the absolute absorption frequencies of nuclei is difficult. A standard internal reference compound is added, and the frequency difference between the reference signal and sample signals is measured.

The internal reference compound generally used in NMR spectroscopy is tetramethylsilane (TMS). TMS is preferred because it is chemically inert, soluble in NMR solvents, and easily removable. Also, the highly shielded methyl protons in TMS yield an intense signal at a lower frequency than most other organic molecules. Because of these advantages, TMS is used as a primary reference in proton, carbon, and silicon NMR spectroscopy. If a suitably inert reference compound is not available, the reference is kept in a capillary tube within the NMR tube and called an external reference.

In addition, deuterated NMR solvents such as CDCl3, D2O, and (CD3)2SO contain residual protons whose signal can be used as a secondary reference. Furthermore, the signal from the deuterium itself can be used to monitor the instrument's magnetic field by a technique called locking. During locking, the deuterium signal is constantly compared to a reference frequency and adjusted if there is any variation. 

 Core: Analytical Chemistry

Complexation Equilibria: Factors Influencing Stability of Complexes

JoVE 14745

In complexation reactions, metal cations are the electron pair acceptors, and the ligands are the electron pair donors. The stability of the metal complexes depends primarily on the complexing ability of the central metal ion and the nature of the ligands. Generally, the complexing ability of the metal ion depends on the size and charge of the ion. As the metal ion size increases, the stability of the metal complexes decreases, provided that the valency of the metal ion and the ligands remain the same. On the other hand, as the charge on the metal ion increases, the stability of the complexes increases.

Metal ions are categorized into classes A and B based on their preferred ligands, specifically, ligands described in terms of soft or hard acids and bases. Soft bases have easily oxidizable donor atoms with high polarizability and low electronegativity whereas hard bases have donor atoms that show low polarizability and high electronegativity. Smaller, more highly charged metal ions with unsaturated outer orbitals are categorized as Class A and prefer to bind to hard bases. In contrast, larger, less charged or uncharged metal ions with filled outer orbitals are categorized as Class B and prefer complexation with soft bases. In terms of ligand effects on the stability of complexes, smaller, more highly charged ligands tend to form more stable complexes, whereas larger ligands tend to form less stable complexes because of the effects of steric repulsion. Ligands with vacant p or d orbitals can form pi bonds with certain metal ions, which in turn exhibits back bonding behavior toward the ligand, both of which serve to stabilize the complex further.

 Core: Analytical Chemistry

Fascicle Arrangement in Skeletal Muscles

JoVE 14861

Fascicles are bundles of muscle fibers in a skeletal muscle. Muscle fascicle arrangement is directly associated with the power and range of motion of various muscles. The configuration of these fascicles can vary, leading to different functional outcomes.

The four primary types of muscle based on fascicle arrangement are:

  1. Parallel Muscles: In this type, fascicles run parallel to the long axis of the muscle. Examples include the biceps brachii and the rectus abdominis. When parallel muscles contract, they shorten in length and increase in diameter, enabling movement. They can contract until they have shortened by about 30% in length. The tension developed during contraction depends on the muscle's total number of myofibrils.
  2. Convergent Muscles: These muscles have fascicles that extend over a broad area but converge on a common attachment site. Examples include the pectoralis muscles. Convergent muscles can adapt to different activities as stimulating different muscle regions can change the direction they pull. However, when the entire muscle contracts, the muscle fibers do not pull as hard on the attachment site as a parallel muscle of the same size.
  3. Pennate Muscles: In pennate muscles, the fascicles attach obliquely to a central tendon that runs the length of the muscle, resembling the pattern of feathers along a quill. This arrangement allows for a higher density of muscle fibers in a given area than parallel muscles. The increased fiber density in pennate muscles translates to greater force production, albeit at the cost of reduced range of motion and speed. Pennate muscles are further categorized based on the exact arrangement of their fascicles: unipennate muscles have fascicles attaching to one side of the tendon, bipennate have them attaching on both sides, and multipennate feature a complex, branching network of fascicles and tendons. This intricate architecture of pennate muscles is a perfect adaptation for tasks requiring strong, forceful contractions.
  4. Circular Muscles: Also known as sphincters, circular muscles have fascicles arranged concentrically around an opening. This arrangement forms a ring-like structure that can constrict or relax to control the passage of substances through the body's various pathways. When the circular muscles contract, the orifice narrows or closes, and when they relax, the opening widens, allowing passage. This mechanism is essential in many physiological processes, such as controlling food movement through the gastrointestinal tract or regulating blood flow through blood vessels. The iris of the eye is another classic example of a circular muscle, where its contraction and relaxation change the size of the pupil, thereby controlling the amount of light that enters the eye.

 Core: Anatomy and Physiology

Muscles that Move the Thigh

JoVE 14879

The thigh's motion is primarily governed by muscles originating in the pelvic girdle and inserted into the femur. One crucial muscle, the iliopsoas, is a combination of the psoas major and the iliacus muscles, sharing a common insertion point on the lesser trochanter of the femur.

Three other significant muscles are the gluteus maximus, gluteus medius, and gluteus minimus. The gluteus maximus originates from the posterior surface of the ilium, sacrum, and coccyx, and the thoracolumbar fascia and inserts on the iliotibial tract and the gluteal tuberosity of the femur. It primarily acts as the chief extensor of the femur. Conversely, in its reverse muscle action, it powerfully extends the torso at the hip joint.

The gluteus medius and minimus originate anterior to the origin of the gluteus maximus and insert on the greater trochanter of the femur. The gluteus medius muscle functions as a potent abductor of the femur at the hip joint and is a common site for intramuscular injection because of its accessible location. Deep to the gluteus maximus muscle are six muscles — the piriformis, obturator internus, obturator externus, superior gemellus, inferior gemellus, and quadratus femoris, which function as lateral rotators of the femur at the hip joint.

Additionally, the tensor fasciae latae muscle is located on the thigh's lateral surface, inserted into the lateral condyle of the tibia. This muscle, along with its tendons and gluteus maximus muscles, helps form the iliotibial tract. As the name suggests, this muscle is enveloped by the fascia lata, a dense connective tissue encircling the entire thigh.

Lastly, on the medial aspect of the thigh are three unique muscles — the adductor longus, adductor brevis, and adductor magnus. These muscles originate from the pubic bone and insert into the femur, adducting the thigh and demonstrating the unique ability to rotate the thigh both medially and laterally. The adductor longus also flexes the thigh while the adductor magnus extends it. The pectineus muscle also plays a part in adducting and flexing the femur at the hip joint.

 Core: Anatomy and Physiology

Cerebellum: Anatomical Regions

JoVE 14912

The cerebellum, also known as the "little brain," is located in the posterior cranial fossa, inferior to the tentorium cerebelli and dorsal to the brainstem. It plays a significant role in motor control, coordination, and proprioception.

Cerebellar Structure

Externally, the cerebellum features a highly convoluted surface with numerous folia (narrow ridges) separated by shallow sulci (grooves). The cerebellum is divided into two hemispheres by a thin median structure known as the vermis. The vermis primarily coordinates axial and proximal limb musculature. The hemispherical division also delineates the cerebellum's internal organization into distinct lobes, each with specific functions. The anterior and posterior lobes, which are demarcated by the primary fissure, are involved in regulating muscular movements and coordination. The flocculonodular lobe, located inferiorly and set apart by the posterolateral fissure, plays a pivotal role in maintaining balance and equilibrium.

The outer layer of the cerebellum is called the cerebellar cortex and is composed of gray matter. The cerebellar cortex contains three layers — the molecular layer, the Purkinje cell layer, and the granular layer. The molecular layer consists of stellate and basket cells, which provide inhibitory input to Purkinje cells. The Purkinje cell layer contains large, flask-shaped neurons with elaborate branching dendritic trees that relay motor impulses. These dendrites receive inputs from two main sources — parallel fibers, which originate from granule cells in the cerebellar cortex, and climbing fibers, which originate from neurons in the inferior olivary nucleus of the brainstem. Finally, the granular layer is densely packed with small, excitatory granule cells, which receive input from mossy fibers and give rise to parallel fibers that synapse onto Purkinje cells.

Beneath this layer of folia lies the white matter, known as the arbor vitae, due to its tree-like appearance. It consists of myelinated axons coursing through the cerebellar cortex, connecting different regions within the cerebellum and facilitating communication between them. Embedded within the white matter are four pairs of deep cerebellar nuclei — the fastigial, globose, emboliform, and dentate nuclei. These nuclei receive input from the cerebellar cortex and send output to various brainstem nuclei and the thalamus.

Connectivity with the Brainstem

The cerebellum connects with the brainstem through three pairs of bundled white matter tracts known as the cerebellar peduncles.

  • • Superior Peduncles: These form the connection between the cerebellum and the midbrain, diencephalon, and cerebrum. They function to integrate motor commands with sensory inputs.
  • • Middle Peduncles: The middle peduncles are the largest of the three and primarily transport motor information from the pons to the cerebellum. They are vital for planning and executing movements.
  • • Inferior Peduncles: These are essential for conveying sensory information from the medulla to the cerebellum. They aid in adjusting motor actions based on proprioceptive feedback.

 Core: Anatomy and Physiology

Cerebrospinal Fluid

JoVE 14928

Cerebrospinal fluid (CSF) is a colorless liquid that flows around the brain and the spinal cord, playing a vital role in the protection, support, and overall function of the central nervous system (CNS). CSF production, circulation, and absorption are tightly regulated processes essential for the brain and spinal cord to function properly.

CSF Production

CSF is produced mainly in the choroid plexus, a network of capillaries and ependymal cells located within the ventricular system of the brain. The total volume of CSF in an adult human is approximately 125-150 mL. About 500 mL of CSF is produced daily, meaning the entire CSF volume is replaced roughly every 5 to 6 hours. The composition of CSF is primarily water (about 99%), with the remaining 1% consisting of proteins, glucose, electrolytes, and other substances. The composition of CSF differs from that of blood plasma, as it has fewer proteins, a lower concentration of glucose, and a higher concentration of chloride ions.

The production of CSF by the choroid plexus involves active secretion, filtration, and selective diffusion. Ependymal cells transport ions, glucose, and other substances from the blood into the ventricular system. They also actively secrete specific molecules and maintain the unique composition of the fluid. The blood-brain barrier (BBB), formed by the tight junctions between endothelial cells of the capillaries and ependymal cells, prevents the passage of many large molecules, pathogens, and toxins from the bloodstream into the CSF.

CSF Circulation

CSF circulates through a network of four interconnected cavities in the cerebrum - two lateral ventricles, the third ventricle, and the fourth ventricle. It flows from the lateral ventricles into the third ventricle through the interventricular foramina. From there, it flows through the cerebral aqueduct into the fourth ventricle. Finally, CSF exits the ventricular system through the median aperture and the two lateral apertures, entering the subarachnoid space of the brain and spinal meninges.

CSF is ultimately absorbed back into the bloodstream via arachnoid granulations, specialized structures that protrude into the venous sinuses. Several factors, including arterial pulsations, respiratory pacing, postural changes, and intracranial pressure, influence the flow of CSF.

Functions of CSF

CSF serves several essential functions in the human body. It acts as a cushion, absorbing mechanical shocks and distributing pressure evenly around the brain and spinal cord. This property helps protect these delicate structures from injury. The relatively heavy brain is suspended in CSF, reducing its effective weight and preventing it from compressing the base of the skull or damaging the underlying neural tissues. Additionally, CSF maintains the optimal chemical environment for neuronal function by regulating the concentrations of ions, nutrients, and waste products. It also helps maintain a stable pH and temperature within the CNS. Finally, CSF serves as a medium for distributing immune cells and antibodies, providing a first line of defense against potential infections and other threats to the CNS.

 Core: Anatomy and Physiology

Autonomic Nervous System

JoVE 14947

The autonomic nervous system (ANS) is a critical component of the peripheral nervous system, primarily responsible for regulating involuntary bodily functions and maintaining homeostasis. It functions in tandem with the central nervous system (CNS) to seamlessly coordinate various physiological processes without the need for conscious control.

The ANS comprises two main divisions: the sympathetic and parasympathetic divisions. These divisions function antagonistically to maintain a dynamic balance, responding adaptively to different physiological demands.

The Sympathetic Division:

The sympathetic division orchestrates the "fight-or-flight" response, preparing the body for action during stress or danger. Key physiological changes triggered by this division include an increased heart rate, dilation of blood vessels, and a heightened release of stress hormones, notably adrenaline and norepinephrine. Norepinephrine, in particular, plays a significant role in activating this response by acting as a neurotransmitter at the synaptic level and as a hormone in the bloodstream, ensuring a comprehensive and sustained reaction to stress.

The Parasympathetic Division:

On the other hand, the parasympathetic division promotes the "rest-and-digest" response. This division aids in relaxation, digestion, and conservation of energy. It counterbalances sympathetic activity by slowing down the heart rate, facilitating decreased respiratory rates, and enhancing digestive processes. The parasympathetic division is pivotal in helping the body to recover, conserve energy, and maintain a state of calm.

Autonomic Tone and Hypothalamic Regulation:

Autonomic tone refers to the prevailing influence exerted by both sympathetic and parasympathetic divisions, crucial for the optimal functioning of various organ systems. The hypothalamus is pivotal in regulating these autonomic responses and orchestrates various autonomic functions through a sophisticated network of neurons. It integrates multiple signals and modulates the balance between sympathetic and parasympathetic outputs to ensure homeostasis.

The interplay between these two divisions is essential in maintaining homeostasis. For instance, during high-stress situations, the sympathetic division dilates pupils for better vision, increases heart rate, widens airways, and suppresses digestion. In contrast, once the crisis subsides, the parasympathetic division restores normalcy by constricting pupils, slowing the heart rate, normalizing respiratory rates, and promoting digestive activities.

 Core: Anatomy and Physiology

Physiology of Smell and Olfactory Pathway

JoVE 14966

Humans detect odors with the help of specialized cells located in the upper part of the nasal cavity, called olfactory receptor neurons (ORNs). ORNs possess hair-like structures called cilia, which are receptive to sensations from the inhaled air. When an odorant molecule binds to a specific receptor on the cell of the cilia, it leads to a series of events that ultimately cause the ORN to send electrical signals to the olfactory bulb in the brain through the olfactory nerves.

The olfactory bulb, located in the front part of the brain, is responsible for processing and recognizing smells. Upon receiving the signals from the ORNs, the olfactory bulb sends information to other parts of the brain, including the amygdala (associated with emotions) and the hippocampus (associated with memory). Integrating smell with other senses helps us better perceive our environment.

The human olfactory system can detect thousands of odors, each with a unique chemical structure. Interestingly, there are no separate receptors for each odorant molecule. Instead, each receptor can detect multiple odors, and the brain interprets the mix of activated receptors to identify the specific odor. Furthermore, the ability to identify smells is influenced by personal experience and cultural factors. We may associate certain smells with particular memories or emotions, leading to a subjective odor perception.

Once an odorant molecule binds to a receptor, it activates a G protein that activates an enzyme called adenylate cyclase. Adenylate cyclase produces a molecule called cyclic adenosine monophosphate (cAMP). The cAMP molecules bind to and open ion channels, allowing positively charged ions like sodium (Na+) and calcium (Ca2+) to flow into the cell. The influx of positively charged ions generates an electrical signal, which travels down the length of the sensory neuron and is transmitted to the olfactory bulb through the olfactory nerve. The signals are integrated and processed in the olfactory bulb, allowing the brain to recognize and distinguish different odors. The olfactory information is sent to other brain parts, including the amygdala (associated with emotions) and the hippocampus (associated with memory).

The olfactory pathway in humans involves inhaling odor molecules into the nose, binding them to specialized receptor cells in the olfactory epithelium. From there, signals are sent to the olfactory bulb, a structure located at the base of the forebrain. The signals are then relayed to two nearby brain regions: the primary and secondary olfactory cortex. The primary olfactory cortex recognizes odors and associates them with memories or emotional responses. In contrast, the secondary olfactory cortex processes sensory information about odors’ intensity, directionality, and duration. Additionally, recent research has shown that some neural pathways from the primary olfactory cortex may even directly connect to other parts of the brain involved in emotion and behavior. This indicates that the sense of smell plays a much more significant role in behavior and emotion than was once thought.

The primary olfactory cortex and areas of the brain responsible for memory are also believed to be involved in pheromone detection. Pheromones are chemical signals secreted by animals (including humans) that influence the behavior or physiology of other members of the same species. In humans, pheromones have been linked to sexual attraction, although this connection is still poorly understood. Recent studies have suggested some components of human sweat may act as pheromones and could be used to communicate emotions or even influence mood. Further research into the role of the olfactory system in behavior and emotion is needed to understand its effects fully.

 Core: Anatomy and Physiology

Synthesis and Regulation of Thyroid Hormones

JoVE 14982

Low blood levels of the thyroid hormones — triiodothyronine (T3) and thyroxine (T4) — signal the hypothalamus to release the thyrotropin-releasing hormone (TRH). TRH then reaches the pituitary gland and stimulates the release of thyroid-stimulating hormone(TSH) into the bloodstream.

Upon reaching the thyroid gland, TSH stimulates the follicular cells' active uptake of iodide ions from the blood. The ions diffuse to the apical surface of the cells and are oxidized to iodine. The iodine is then added to tyrosine residues of the thyroglobulin protein stored in the lumen colloid, forming monoiodotyrosine (T1) and diiodotyrosine (T2). The coupling of one molecule of T1 with T2 forms a T3 molecule, whereas two T2 molecules combine together to create T4.

Finally, endocytosis and enzymatic cleavage release the T3 and T4 in the follicular cells. These hormones then cross into the blood vessels for transport to target tissues. An increase in blood levels of T3 and T4 temporarily inhibits the production of TRH and TSH, regulating the level of thyroid hormones in the body.

 Core: Anatomy and Physiology

Circuit Terminology

JoVE 15067

An electrical network is a system composed of interconnected elements, such as resistors, capacitors, inductors, and voltage or current sources. Unlike a circuit, an electrical network does not necessarily form a closed path. In other words, while all circuits can be considered networks due to their interconnected nature, not every network qualifies as a circuit.

A circuit, on the other hand, is also an interconnected system of electrical elements but must contain one or more closed paths. These closed paths, or loops, allow current to flow continuously, making the operation of the circuit possible.

The structure of a circuit can be broken down further into several key components: branches, nodes, loops, and meshes. Understanding these components is crucial for analyzing and designing circuits.

A branch represents a single element of the circuit, like a resistor, a voltage source, or a current source, connected by two terminals. For instance, a circuit that includes a voltage source, a current source, and two resistors would consist of four branches.

Nodes are points in the circuit where two or more branches intersect. They are typically represented by a dot in circuit diagrams. If a short circuit connects multiple nodes together, they effectively become a single node.

A loop in a circuit is a closed path that passes through several nodes without traversing any node twice. Loops can be independent or dependent. An independent loop has at least one branch that is not shared with any other loop.

A mesh, on the other hand, is a special type of loop. It is a closed path that does not contain any other loops within it. Meshes are fundamental to certain types of circuit analysis, such as mesh analysis, which systematically solves for current values.

Lastly, all these elements are tied together by the fundamental theorem of network topology, which states that for a network containing 'b' branches, 'n' nodes, and 'l' independent loops, these quantities are related by a specific formula, given by

Equation1

Understanding this relationship is crucial for the analysis and design of electrical networks and circuits.

 Core: Electrical Engineering

Instrumentation Amplifier

JoVE 15083

An electrocardiography (ECG) machine is an essential piece of medical equipment used to monitor the electrical activity of the heart. It operates by detecting small electrical changes on the skin that result from the depolarization of the heart muscle during each heartbeat. However, these signals are in the microvolt range and can be easily overwhelmed by noise or interference.

To overcome this challenge, an ECG machine utilizes an instrumentation amplifier. This specialized amplifier is designed to amplify the ECG waveform while minimizing the impact of common-mode noise signals, such as electrical interference from other devices.

The instrumentation amplifier is a type of difference amplifier. Its circuitry includes six resistors, three terminals, and an external resistor connected between the gain set terminals. This amplifier functions by amplifying the small differences between input signals while rejecting signals common to both inputs.

The output voltage of the amplifier is determined by the product of the voltage gain and the input voltage difference. The gain can be easily adjusted by varying the value of the external resistor, providing flexibility and control over the amplification process.

Instrumentation amplifiers exhibit several key characteristics that make them ideal for use in ECG machines. They have high input impedance, low output impedance, high gain stability, and a high common mode rejection ratio (CMRR).

The high input impedance prevents signal distortion by avoiding signal loading, ensuring the integrity of the electrical signals from the heart. The high CMRR enables the amplifier to effectively reject noises picked up by the electrode leads, such as electrical interference from other equipment, thereby improving the quality of the ECG waveform.

Beyond their use in ECG machines, instrumentation amplifiers are commonly employed in a variety of biomedical instrumentation and telecommunication applications due to their unique properties.

 Core: Electrical Engineering

Series RLC Circuit with Source

JoVE 15100

Consider the operation of an automobile ignition system, a crucial component responsible for generating a spark by producing high voltage from the battery. This system can be described as a simple series RLC circuit, allowing for an in-depth analysis of its complete response.

In this context, the input DC voltage serves as a forcing step function, resulting in a forced step response that mirrors the characteristics of the input. Applying Kirchhoff's voltage law to the circuit yields a second-order differential equation. Remarkably, this equation strongly resembles the second-order differential equation characterizing a source-free RLC circuit. This similarity underscores that the presence of the DC source does not alter the fundamental form of the equations.

Equation1

The complete solution to this equation comprises both transient and steady-state responses.

Equation2

The transient response, which diminishes over time, aligns with the solution for source-free circuits and encompasses scenarios involving overdamped, critically damped, and underdamped behaviors. On the other hand, the steady-state response corresponds to the final value of the capacitor voltage, which equals the source voltage. The constants involved in these responses can be determined from the initial conditions of the circuit.

 Core: Electrical Engineering

Linear Circuits

JoVE 15170

A linear circuit is characterized by its output having a direct proportionality to its input, adhering to the linearity property, which encompasses the principles of homogeneity (scaling) and additivity. Homogeneity dictates that when the input, also referred to as the excitation, is multiplied by a constant factor, the output, known as the response, is correspondingly scaled by the same constant factor. For instance, if the current is multiplied by a constant 'k,' the voltage likewise experiences an increase of 'k' times.

The additivity property stipulates that the response to a sum of inputs equals the sum of responses to each individual input applied separately. In essence, this property enforces that the circuit's behavior remains consistent even when multiple inputs are combined.

Notably, a resistor is classified as a linear element because it satisfies both the homogeneity and additivity properties within its voltage-current relationship. Generally, a circuit is considered linear if and only if it demonstrates both additivity and homogeneity characteristics. Such linear circuits exclusively comprise linear elements, linear dependent sources, and independent sources.

Conversely, the expression for power, which is defined as the ratio of the square of voltage to resistance, constitutes a quadratic function and, therefore, falls under the category of nonlinearity within the context of circuit analysis.

 Core: Electrical Engineering

Leveling Effect

JoVE 17365

In acid-base chemistry, the leveling effect refers to the limitation imposed by the solvent on the strength of acids and bases in solution. When a base stronger than the solvent's conjugate base is used, it deprotonates the solvent until the base is entirely consumed, making it ineffective against weaker acids. Conversely, an acid stronger than the solvent's conjugate acid protonates the solvent until the acid is depleted, rendering it ineffective against weaker bases. Essentially, the solvent neutralizes stronger acids and bases, preventing them from reacting as intended with other compounds.

For example, in water (an aqueous solution), a strong base like the amide ion deprotonates water, predominantly forming hydroxide ions and leaving few amide ions. This prevents the amide ions from effectively deprotonating compounds like acetylene, which have a higher pKa than water. However, if a more basic solvent like ammonia is used, the amide ions can successfully deprotonate acetylene, facilitating the desired reaction.

In summary, the solvent's acidity limits the effectiveness of strong bases, and its basicity limits the effectiveness of strong acids. The chosen solvent must, therefore facilitate the acid-base reaction without undergoing a significant reaction itself.

 Core: Analytical Chemistry

Introducing Experimental Agents into the Mouse

JoVE 5161

Many investigations performed in mice (Mus musculus) require the administration of an experimental agent to the animal. For example, it may be of interest to test the efficacy of a specific therapy, to induce a pathologic condition, or to administer anesthesia or palliative care. In order to ensure safe and efficient delivery, it is important to consider a variety of factors prior to the administration of the treatment.

This video, which reviews agent administration in the mouse, begins by highlighting properties to consider, such as viscosity, dose, and palatability, when planning the administration of an experimental agent. The subsequent discussion focuses on injection methods, including delineation of the structural components of the syringe and needle, how to interpret needle gauge, and safe mouse restraint methods for common injection sites. Detailed instructions are provided for performing subcutaneous (SC/SubQ), intraperitoneal (IP), and tail vein (IV) injections in mice. Furthermore, applications of these techniques as well as alternative administration routes are discussed.

 Biology II

Primary Neuronal Cultures

JoVE 5214

The complexity of the brain often requires neuroscientists to use a simpler system for experimental manipulations and observations. One powerful approach is to generate a primary culture by dissecting nervous system tissue, dissociating it into single cells, and growing those cells in vitro. Primary cultures make neurons and glia easily accessible to the experimental tools required for techniques like genetic manipulation and time-lapse imaging. Furthermore, these cultures represent a highly controllable environment in which to study complex phenomena such as cell-cell interactions.

This video provides an overview of the major steps in producing primary neuronal cultures, which include selecting and dissecting the tissue of interest, mechanically and chemically breaking down the tissue to produce a single cell suspension, plating the cells, and maintaining the cultures in the appropriate media. Several example experiments are also presented to show how cultured cells can be used to investigate protein trafficking, morphological changes, and electrophysiology in living neurons.

 Neuroscience

Transplantation Studies

JoVE 5336

Many developmental biologists are interested in the molecular signals and cellular interactions that induce a group of cells to develop into a particular tissue. To investigate this, scientists can use a classic technique known as transplantation, which involves tissue from a donor embryo being excised and grafted into a host embryo. By observing how transplanted tissues develop in host environments, scientists have started to dissect the molecular pathways underlying development.

In this video, we first look at the role of cellular interactions in development, and move on to a basic transplantation protocol. Finally, some specific developmental studies utilizing this technique are discussed, which examine the effect of tissue transplantation on the fate of donor and host tissue.

 Developmental Biology

Positive Reinforcement Studies

JoVE 5426

Researchers study learning of a behavior through the use of operant conditioning. This type of learning involves associating the behavior with a consequence, which is a reward or punishment. If the consequence is a reward, it leads to reinforcement of the desired behavior. One type of reinforcement approach is positive reinforcement, where the behavior is rewarded with an artificial, natural, or social reinforcer. Studies using positive reinforcement as a tool can help tease out important details about neurological functioning associated with different behaviors.

This video reviews the concepts behind reinforcement studies by using an example of a man training a dog to sit. Following this, we look at a generalized procedure of positive reinforcement commonly used by behavioral researchers. This involves, training rodents to perform a behavior (lever press) to get a reward (food). Lastly, specific applications demonstrate how scientists use positive reinforcement to understand behavior.

 Behavioral Science

Genetic Screens

JoVE 5542

Genetic screens are critical tools for defining gene function and understanding gene interactions. Screens typically involve mutating genes and then assessing the affected organisms for phenotypes of interest. The process can be “forward”, where mutations are generated randomly to identify unknown genes responsible for the phenotypes, or it can be “reverse”, where specific genes are targeted for mutation to observe what phenotypes are produced.

Here, JoVE reviews various types of genetic screens, including those that depend on either loss-of-function or gain-of-function mutations, which respectively decrease or increase the activity of genes. We then explore general protocols for forward and reverse screens in a popular model organism, the nematode worm. Finally, we highlight how screens are applied in research today, for example to better understand gene interactions that may contribute to neurodegenerative diseases.

 Genetics

The ATP Bioluminescence Assay

JoVE 5653

In fireflies, the luciferase enzyme converts a compound called luciferin into oxyluciferin, and produces light or “luminescence” as a result. This reaction requires energy derived from ATP in order to proceed, so researchers have exploited the luciferase-luciferin interaction to gauge ATP levels in cells. Given ATP’s role as the cell’s currency of energy, the ATP bioluminescence assay can provide insight into cellular metabolism and overall cell health.

In this video, JoVE discusses cellular respiration, specifically reviewing how glucose metabolism results in ATP production. This is followed by principles behind the ATP bioluminescence assay and a generalized protocol for this technique. Finally, a survey of how researchers are currently using the ATP bioluminescence assay to evaluate cell viability in a variety of experimental conditions.

 Cell Biology

Förster Resonance Energy Transfer (FRET)

JoVE 5696

Förster resonance energy transfer (FRET) is a phenomenon used to investigate close-range biochemical interactions. In FRET, a donor photoluminescent molecule can non-radiatively transfer energy to an acceptor molecule if their respective emission and absorbance spectra overlap. The amount of energy transferred—and consequently the overall emission of sample—depends on the proximity of an acceptor-donor pair of photoluminescent molecules. FRET analysis is combined with other biochemistry techniques to obtain detailed information of biomolecular structures and interactions from this “spectroscopic ruler.”

This video covers the principles and concepts of FRET analysis. The procedure focuses on preparing samples for FRET and ways to present and interpret data. Finally, the applications include monitoring conformational and cellular processes by labeling parts of a cell or protein, monitoring enzyme reactions that alter protein structures, and using FRET to monitor aggregation of monomers expressed by cells.

Förster Resonance Energy Transfer, or FRET, is a non-radiative transfer of energy between light-emitting molecules, and is often used to investigate close-range biochemical interactions. FRET only occurs when fluorescent molecules are spaced within 10 nm of each other. FRET analysis can be combined with other techniques to obtain detailed structural information. This video will introduce the underlying principles of FRET, summarize a protocol and data presentation, and discuss some biochemical applications.

A photoluminescent molecule such as a fluorophore is excited by absorbing electromagnetic radiation at a wavelength in its absorption spectrum. As it relaxes, it emits light at a wavelength within its emission spectrum. For more information about fluorescence, see JoVE's video on fluorescence microscopy. Different fluorophores absorb and emit light at different wavelengths, which frequently overlap. If the emission spectrum of a fluorophore overlaps significantly with the absorption spectrum of another fluorophore, the “donor” will release a virtual photon, which is absorbed by the “acceptor”. When an excited donor is within 10 nm of an acceptor, energy is transferred from donor to acceptor by dipole-dipole interactions. The release of energy by emission of light from the donor correspondingly decreases. Meanwhile, the excited acceptor emits light at its emission wavelength. The FRET response is evaluated in terms of efficiency, or the percentage of energy released from the donor by FRET rather than by fluorescence or other radiative processes. The efficiency depends strongly on the distance between the donor and acceptor, which allows FRET to act as a 'molecular' or 'spectroscopic' ruler.

In biochemistry, FRET is often used qualitatively to observe conformational changes in molecules by monitoring fluorophores as they move in and out of FRET range of each other. Similarly, cellular functions can be studied with molecules containing a FRET pair. If the labeled molecule is cleaved by enzyme activity, FRET stops and the observed fluorescence wavelength changes.

Now that you understand the principles behind FRET, let's look at an overview of a protocol and a few ways to present and interpret the data.

Prior to the experiment, the biomolecules of interest, typically DNA or proteins, are engineered with fluorescent tags, using molecular biology techniques. Common ways to introduce the modified genetic material into the cells include transfection and electroporation.

Then, the cells are prepared for FRET visualization on a fluorescence microscope. For instance, the molecules may be immobilized on a slide for single-molecule FRET, or samples are loaded into wells for high-throughput screening.

Then, the excitation lasers, microscope, and associated equipment are prepared. (A) FRET experiments often involve powerful lasers; (B) so appropriate PPE and safety procedures should be used. The sample is then placed in the instrument and illuminated with the excitation laser.

For experiments monitoring cell behavior, color images showing differences or changes in emission intensity are used. Donor and acceptor emission intensities are plotted together to track FRET response over time.

FRET data can also be fitted to various functions for more complex analyses. Depending on the experiment, data may be presented in multiple ways to best represent the results, making FRET a flexible experimental tool.

Now that you're familiar with the basics of running and analyzing a FRET experiment, let's look at some applications of FRET in biochemistry research.

FRET can be used to study conformational changes or cellular processes by labeling parts of the protein or cell predicted to move within 10 nm of each other with a FRET pair. For example, protein sensors are prepared by labeling receptors with a pair of fluorophores. The FRET response is monitored live by confocal microscopy. Variation of emission wavelength and intensity indicate conformational changes.

FRET can also be used by preparing molecules with an active FRET pair and observing changes in the response. When the substrate is cleaved, FRET is disrupted, causing an increase in donor emission and a decrease in acceptor emission. The emissions are analyzed to determine contributions by donor, acceptor, and FRET. Once the direct emission factors are calculated for the cyan and yellow fluorescent proteins, the concentration and kinetic parameters of the substrate can be determined.

Cells designed to express monomers containing either of a FRET pair function as 'sensors' for interactions between those monomers. If aggregation of those monomers is induced, a FRET response is observed. This can be used to investigate protein aggregation triggered by 'seeding' of misfolded proteins. Here, cells were transduced with aggregates of the protein of interest, incubated, and analyzed with flow cytometry.

You've just watched JoVE's video on Förster Resonance Energy Transfer, or FRET. This video contained the underlying principles of FRET, preparation and analysis of a FRET experiment, and a few biochemical applications.

Thanks for watching!

 Biochemistry

Optical Biosensing

JoVE 5795

Optical biosensors utilize light to detect the binding of a target molecule. These sensors can utilize a label molecule, which produces a measurable signal such as fluorescence, or these sensors can be label-free and use the changes in optical properties, such as refractive index, to sense for the binding of the target molecule. This video introduces both label and label-free optical biosensors, demonstrates their use in the laboratory, and shows some applications of the technology.

 Bioengineering

Conformations of Cycloalkanes

JoVE 11711

Adolf von Baeyer attempted to explain the instabilities of small and large cycloalkane rings using the concept of angle strain — the strain caused by the deviation of bond angles from the ideal 109.5° tetrahedral value for sp3  hybridized carbons. However, while cyclopropane and cyclobutane are strained, as expected from their highly compressed bond angles, cyclopentane is more strained than predicted, and cyclohexane is virtually strain-free. Hence, Baeyer’s theory that was based on the assumption that all cycloalkanes are flat was wrong, and, in reality, most cycloalkanes adopt a non-planar structure.

Cyclopropane, the three-carbon cyclic alkane, has the highest angle strain since its planar structure is highly compressed, deviating by 49.5° from the ideal value. Additionally, cyclopropane has a torsional strain due to the eclipsing interaction between six C-H bonds. Hence, cyclopropane has an overall ring strain of 116 kJ/mol. Unlike cyclopropane, which is planar, cyclobutane takes up a more stable, folded non-planar conformation. Folding causes the angle strain to be slightly elevated compared to the hypothetical planar cyclobutane, but the torsional strain from the ten eclipsing hydrogens is greatly relieved. Cyclobutane has an overall strain of 110 kJ/mol. Cyclopentane also adopts a non-planar conformation known as envelope conformation. Compared to the hypothetical planar form of cyclopentane, the envelope form has its bond angles slightly reduced, which marginally increases the angle strain. However, it significantly alleviates the torsional strain from ten eclipsing C-H bonds. Hence, the overall strain in cyclopentane is 27 kJ/mol.

 Core: Organic Chemistry

Autoxidation of Ethers to Peroxides and Hydroperoxides

JoVE 11763

Ethers represent a class of chemical compounds that become more dangerous with prolonged storage because they tend to form explosive peroxides when standing in the air. Autoxidation is the spontaneous oxidation of a compound in air. In the presence of oxygen, ethers slowly oxidize to form hydroperoxides and dialkyl peroxides.

Figure1

If concentrated or heated, these peroxides may explode. Hence, ethers should be obtained in small quantities, kept in tightly sealed containers, and used promptly to prevent such explosions. Autoxidation of ethers proceeds by a free-radical chain reaction consisting of a series of steps—initiation, propagation, and termination in repetitive cycles. Each of these steps forms intermediate products called chain carriers that regenerate in each step. Such a reaction will continue as long as the chain carriers persist. Hydroperoxides and peroxides can be detected by shaking ether samples with an acidified aqueous 10% solution of potassium iodide, thereby liberating iodine which gives the yellow color to the solution.

 Core: Organic Chemistry

Oxidation of Alkenes: Syn Dihydroxylation with Potassium Permanganate

JoVE 11782

Alkenes can be dihydroxylated using potassium permanganate.  The method encompasses the reaction of an alkene with a cold, dilute solution of potassium permanganate under basic conditions to form a cis-diol along with a brown precipitate of manganese dioxide.

Figure1

The mechanism begins with the syn addition of a permanganate ion (MnO4) across the same side of the alkene π bond, forming a cyclic manganate ester intermediate. Next, the hydrolysis of the cyclic ester with water gives a cis-diol with the retention of stereochemistry at the newly formed C–O bonds.

Figure2

Potassium permanganate is inexpensive and safer compared to osmium tetroxide. However, its strong oxidizing nature leads to over-oxidation of the diol, thereby giving poor yields.

Syn Dihydroxylation with Hot Basic Potassium Permanganate

When hot potassium permanganate is used, it oxidatively cleaves the carbon–carbon double bond, forming ketones or acids depending on the nature of the substituents on the alkene. Thus, terminal alkenes are oxidized to form carbon dioxide, while monosubstituted and disubstituted alkenes give carboxylic acids and ketones, respectively.

Figure3

Qualitative Analysis Using Potassium Permanganate

The basic potassium permanganate solution is also known as Baeyer's reagent, which is used in qualitative analysis to determine the presence of olefinic double bonds. During the reaction, the deep purple color of potassium permanganate solution decolorizes with the formation of a brown precipitate of manganese dioxide.

 Core: Organic Chemistry

Alkynes to Carboxylic Acids: Oxidative Cleavage

JoVE 11840

Alkynes undergo oxidative cleavage in the presence of oxidizing reagents like potassium permanganate and ozone. The triple bond — one σ bond and two π bonds — is completely cleaved, and the alkyne is oxidized to carboxylic acids. When warm and basic aqueous potassium permanganate is used as an oxidizing agent, alkynes are first converted to carboxylate salts via an unstable α-diketone intermediate. Further, a mild acid treatment protonates the carboxylate anions generating free carboxylic acid molecules. When an alkyne is subjected to ozonolysis, an ozonide intermediate is formed, which is then oxidatively cleaved through hydrolysis to yield carboxylic acids.

Oxidative cleavage of internal alkynes yields only carboxylic acids, while terminal alkynes generate carbon dioxide, in addition to an acid, irrespective of the oxidizing reagent used. Thus, oxidative cleavage can be used to locate the triple bond in unknown alkynes. The carbonyl groups in the products are a key to determine the position of the oxidatively cleaved triple bond in the reactant: if the identity of the acids is known, the structure of the unknown alkyne can be deduced.

 Core: Organic Chemistry

Acidity and Basicity of Alcohols and Phenols

JoVE 11922

Like water, alcohols are weak acids and bases. This is attributed to the polarization of the O–H bond making the hydrogen partially positive. Moreover, the electron pairs on the oxygen atom of alcohol make it both basic and nucleophilic. Protonation of an alcohol converts hydroxide, a poor leaving group, into water—a good one. The two acid–base equilibria corresponding to ethanol are depicted below.

Figure1

Figure 1. Loss of proton

Figure2

Figure 2. Gain of proton

Methanol (pKa = 15.5) is the only alcohol that is slightly stronger than water (pKa = 15.7). Ethanol (pKa = 15.9), tert-butanol (pKa = 18.0), and others are weaker acids. However, all alcohols are stronger acids than terminal alkynes, and they are much stronger than hydrogen, ammonia, and alkanes.

Figure3

Figure 3. Relative acidity

Since alcohols are weaker acids than water, their conjugate base alkoxide ions are stronger bases than the hydroxide ion. Alcohol can be converted into metal alkoxide using strong bases such as sodium/potassium hydride or sodium/potassium metal, which react violently but controllably with alcohol. When the alkyl substitution is bulky, the alkoxide ion is not solvated enough due to the steric effect, leading to less stabilization. Additionally, destabilization is favored by inductive effects. Consequently, the equilibrium lies predominantly towards alcohol.

Alcohols can also act as a base and accept protons from strong acids. Notably, conjugate bases of compounds with higher pKa than an alcohol will deprotonate that alcohol.

Figure4

Figure 4. Relative basicity

Phenols are more acidic than alcohols. The conjugate base of a phenol is a phenoxide or phenolate ion. Resonance stabilization of the phenoxide ion coupled with the polar effect of the benzene ring enhances the acidity of phenols by eight orders of magnitude (100,000,000 times) over cyclohexanol. Therefore, phenol does not need to be deprotonated with a base as strong as sodium hydride. Instead, it can be deprotonated by hydroxide, unlike an alcohol.

Although phenol is a million-fold higher in acidity than ethanol, it is a hundred thousand-fold less acidic than acetic acid. Their relative acid–base properties can be used to separate each other from a mixture. When an ether solution of a mixture of alcohol and phenol is extracted with dilute sodium hydroxide, phenol gets completely partitioned into the aqueous phase as its sodium salt, while alcohol stays in the ether layer. On the other hand, dilute sodium bicarbonate is used to extract phenol and carboxylic acid from an ether solution of their mixture. Carboxylic acid gets quantitatively converted into its sodium salt and gets extracted from water while phenol remains in the ether phase.

Figure5

Figure 5. Acid–base equilibria of phenol

A phenoxide ion can be stabilized by delocalization of the oxygen's negative charge on the benzene ring. This is reinforced by electron-withdrawing functional groups like nitro, halide, etc. A substantial change in acidity is noted in phenols with an electron-withdrawing substituent, like a nitro group. An ortho- or para-nitro group stabilizes the phenoxide ion by delocalizing the negative charge on its own oxygen atoms. On the other hand, a meta-nitro group, being not directly conjugated to the phenoxide oxygen, stabilizes the phenolate ion to a lesser extent. Therefore, m-nitrophenol (pKa = 8.4) is more acidic than phenol but less acidic than o- or p-nitrophenol (pKa = 7.2). This also explains the extremely high acidity of 2,4-⁠dinitrophenol (pKa = 4.0) and 2,4,6-trinitrophenol (pKa = 0.4).

Table 1. The acidity constants (pKa) of acids (blue) and their conjugate bases.

Compound Acid Conjugate base pKa
Hydrogen chloride Figure6 Figure7 −6.30
Nitrophenol Figure8 Figure9 7.07
Phenol Figure10 Figure11 9.89
m-Cresol Figure12 Figure13 10.1
2,2,2-Trifluoroethanol Figure14 Figure15 12.0
Water Figure16 Figure17 15.7
Ethanol Figure18 Figure19 15.9
Cyclohexanol Figure20 Figure21 16.0

 Core: Organic Chemistry

Microtubules in Signaling

JoVE 12135

The primary cilium, made up of microtubules, acts as antennae on the cell surfaces for relaying external stimuli into the cells. These fine hair-like structures are present, generally one per cell. These are non-motile cilia in a 9+0 microtubules arrangement, where the central pair of microtubules are absent. The primary cilia arise from the basal body embedded in the cell membrane. Intraflagellar transport (IFT) carries requisite proteins from the cytoplasm to the cilium because the primary cilium cannot synthesize proteins. These are resistant to microtubule disassembly-causing drugs like colchicine.

The primary cilium plays a role in several key signaling pathways, among which the calcium-dependent signaling pathway has been most widely studied. However, recent studies have highlighted their roles in calcium-independent pathways like Sonic-hedgehog, Wnt, PDGFR, Notch, etc.

In humans, the primary cilium, although found in almost all cell types, but is most common in epithelial cells. They have a key role in the node of the vertebrate, which is responsible for positioning the organs in the developing embryo. The primary cilium moves in a circular motion to create the left-right symmetry for the correct positioning of the visceral organs. A defect in genes responsible for forming primary cilium results in the sinus invertus disorder, where the typical symmetry of the organs in the embryo is lost. Other diseases associated with the defect in primary cilium include Meckel-Guber syndrome, Bardet-Beedle syndrome, polycystic kidney disorder, and Joubert syndrome.

 Core: Cell Biology

The Intrinsic Apoptotic Pathway

JoVE 12429

Internal cellular stress, such as cellular injury or hypoxia, triggers intrinsic apoptosis. The B-cell lymphoma 2 (Bcl-2) family of proteins are the primary regulators of the intrinsic apoptotic pathway. For example, during DNA damage, checkpoint proteins, such as Ataxia Telangiectasia Mutated (ATM protein) and Checkpoints Factor-2 (Chk2) proteins, are activated. These proteins phosphorylate p53 which further activates pro-apoptotic proteins, such as Bax, Bak, PUMA, and Noxa, and inhibits anti-apoptotic proteins Bcl-2, Bcl-XL, and survivin. This whole chain of events eventually leads to apoptosis of the damaged cells.

Cancer and the Intrinsic Apoptotic Pathway

The intrinsic pathway is inhibited during cancer due to mutations in pro-apoptotic and anti-apoptotic proteins. For example, in melanoma, the overexpressed inhibitor of apoptosis proteins or IAPs causes the inactivation of caspase-9, preventing apoptosis. In most cancer cells, Bcl-2, an anti-apoptotic protein first found in follicular lymphoma, is overexpressed.

Several anti-cancer drugs that can activate the intrinsic apoptotic pathway have been developed. For example, molecules such as Nultlin-2 and MI-219 inhibit the binding of MDM2 to p53, thereby preventing the inactivation of p53 and leading to apoptosis of cancerous cells. Some drug molecules, such as SH122, JP1201, and YM155, inhibit IAPs and help activate caspases, thereby initiating apoptosis in cancer cells.

 Core: Cell Biology

Growth of Cartilage and Bone Tissue

JoVE 12519

Chondrocytes form a temporary cartilaginous model by dividing and secreting a thick gel-like extracellular matrix. Once the chondrocytes undergo programmed cell death, osteoblasts enter the site of the cartilaginous model. The process of replacing the temporary cartilaginous model with bone in an ordered manner is called endochondral ossification. In endochondral ossification, not all of the cartilage is replaced by bone tissue. Some cartilage that performs a protective and supportive function remains intact.

Cartilage protects the ends of long bones. It also maintains the shape of flexible body parts such as the external ears and the larynx. These body parts contain a type of cartilage called elastic cartilage. The second type of cartilage, hyaline cartilage, is present in organs such as the nose, ends of ribs, and the ends of bones that make up the joints of the skeleton. Hyaline cartilage reduces friction between bones and provides flexibility to the joints. The third type of cartilage, fibrous cartilage, is found in tissues such as intervertebral discs and ligaments.

Cartilage can be damaged through wear and tear, such as sports injuries or osteoarthritis. People with damaged cartilage experience pain in joints, swelling, or stiffness. Cartilage is avascular, so nutrients must diffuse into the tissue; therefore, cartilage injuries often take time to heal. Additionally, cartilage is incapable of regeneration in contrast with bones that continuously undergo remodeling.

A mature bone consists of four types of tissues. The hard outer part of the bone is the compact tissue. Below that exists a sponge-like tissue called the cancellous tissue. The tissue at the ends of the bones, protected by cartilage, is called the subchondral tissue. Mechanical stress to a bone, such as a fracture, heals in several stages, which include hematoma formation, fibrocartilaginous callus formation, bony callus formation, and bone remodeling. The last stage, bone remodeling, is the longest and is where endochondral ossification occurs.

 Core: Cell Biology

Elastin is Responsible for Tissue Elasticity

JoVE 13340

Elastic fiber contains the protein elastin along with lesser amounts of other proteins and glycoproteins. The main property of elastin is that it will return to its original shape after being stretched or compressed. Elastic fibers are prominent in elastic tissues found in skin and the elastic ligaments of the vertebral column.

Ligaments and tendons are made of dense regular connective tissue, but in ligaments not all fibers are parallel. Dense regular elastic tissue contains elastin fibers and collagen fibers, allowing the ligament to return to its original length after stretching. The ligaments in the vocal folds and between the vertebrae in the vertebral column are elastic. Dense irregular elastic tissues give arterial walls the strength and the ability to regain original shape after stretching

While older adults are at risk for tendinitis because the elasticity of tendon tissue decreases with age, active people of all ages can develop tendinitis. Young athletes, dancers, and computer operators; anyone who performs the same movements constantly is at risk for tendinitis.

Elastic cartilage contains elastic fibers as well as collagen and proteoglycans. This tissue gives rigid support as well as elasticity. Tug gently at your ear lobes, and notice that the lobes return to their initial shape. The external ear contains elastic cartilage.

This text is adapted from Openstax, Anatomy and physiology 2e, Section 4.3: Connective tissue supports and protects.

 Core: Cell Biology

SDS-PAGE

JoVE 13380

Gel electrophoresis is a method that separates biological macromolecules like nucleic acids or proteins by forcing them to pass through a gel matrix under an electric field.

A variation of gel electrophoresis, termed  polyacrylamide gel electrophoresis (PAGE), is commonly used for separating proteins according to their molecular size by passing them through a polyacrylamide gel. Because of the varying charges associated with amino acid side chains, PAGE can be used to separate intact proteins based on their net charges and size, as seen in native PAGE. Alternatively, proteins can be denatured and coated with a negatively charged detergent called sodium dodecyl sulfate (SDS), masking the native charges and allowing separation based on size only, as seen in SDS-PAGE. Polypeptide chains can migrate in polyacrylamide gel even when coated with SDS. So, with this observation, Ulrich K. Laemmli developed the technique of SDS-PAGE in the year 1970.

In SDS-PAGE, the polyacrylamide gels are formed by the polymerization of acrylamide monomers that are transparent, chemically and biologically inert, and uncharged. These gels  have a controllable pore size determined during gel preparation, where the concentration of acrylamide and bisacrylamide, the cross-linking agent, regulate the gel’s pore size. Acrylamide monomers thus polymerize to form polyacrylamide, wherein ammonium persulfate catalyzes the polymerization reaction. SDS provides a uniform charge to mass ratio for all proteins. Usually, a gram of protein is covered by 1.4 gm of SDS allowing size-driven protein separation.

The sample preparation buffer for SDS-PAGE, besides SDS and β-mercaptoethanol, contains glycerol and bromophenol blue. The density of glycerol helps the sample reach the bottom of the stacking well, preventing it from flowing out from the well into the buffer. Bromophenol blue functions as the tracking dye and indicates the proteins’ progress in the gel.

Optimum pH for each buffer solution is crucial as it determines ion concentration in the buffer required for protein movement under voltage application. The ionic strength and pH of the buffer used for running the gel (pH 8.3) is different from the buffers used to create  the stacking gel (pH 6.8) and the resolving gel (pH 8.3). This pH difference between the stacking and resolving gel ensures that the low ionic strength of the stacking gel offers high electrical resistance, allowing the proteins to stack and then separate once they enter the resolving gel.

SDS-PAGE has varied applications, including the estimation of protein size, purity, and even peptide mapping. The major limitation of this technique is that it cannot obtain information about the enzyme activity, cofactor, and protein binding interactions. It is challenging to analyze highly acidic or basic proteins using SDS-PAGE.

 Core: Cell Biology

Total Internal Reflection Fluorescence Microscopy

JoVE 13396

Total internal reflection fluorescence microscopy or TIRF is an advanced microscopic technique used to visualize fluorophores in samples close to a solid surface with a higher refractive index, such as a glass coverslip. TIRF only allows fluorophores in proximity to the solid surface to be excited. When light from a medium with a lower refractive index (such as air) hits the glass coverslip at a critical angle, the light undergoes total internal reflection stead of passing through the glass. This happens as the sample has a lower refractive index than the coverslip and does not allow the light to enter within. The light is reflected from the interface and forms an electromagnetic field emitting short-length evanescent waves. These waves only excite the fluorophore near the surface as they can move only about 100 to 200 nm deep within the cell before dying out.

There are two types of TIRF; prism-based and objective-based. In prism-based TIRF microscopy, a prism is placed on the coverslip surface that directs the evanescent wave to the sample. In objective-based TIRF microscopy, there is no prism at the interphase; the objective is the same as the light source that helps create the evanescent wave.

TIRF has several advantages over traditional fluorescence and confocal microscopes; it prevents the illumination of background fluorophores. It helps in studying the structures close to the cell surface. It reduces the blurring effect and does not allow out-of-focus light to interfere with the image. As the samples are not directly exposed to an intense light beam, the photobleaching is minimum, and the cells are less exposed to phototoxicity.

 Core: Cell Biology

Introduction to Nuclear Reprogramming

JoVE 13476

Nuclear reprogramming is the process of switching gene expression of one cell type to that of another cell type, usually from a differentiated cell state to an undifferentiated cell state. Differentiation occurs during processes such as development and morphogenesis, tissue regeneration, and malignancy. Cells can also be artificially induced to reprogram their gene expression by techniques such as nuclear transfer, induced pluripotency, and cell fusion. Such techniques have many applications in the fields of molecular medicine, biotechnology, and oncology.

A Brief History

In 1952, Briggs and King transplanted a nucleus from a frog blastula into an enucleated frog egg. The resulting embryo developed into a normal tadpole, demonstrating nuclear reprogramming. A few years later, Gurdon and colleagues performed a similar experiment using the nuclei of intestinal epithelia from adult frogs. Despite the nucleus being fully differentiated and specialized, the transplanted eggs developed into swimming tadpoles. This experiment led to the first use of the word "clone" in reference to animals.

Dolly - The First Clone

Dolly was the first animal cloned from an adult somatic cell. Nuclei from the mammary gland cells of an adult sheep were harvested and transferred to an enucleated sheep ovum in a procedure called somatic cell nuclear transfer (SCNT). The egg was then transplanted into a female sheep, who carried it to term and gave birth to Dolly. This successful cloning spawned attempts to clone other mammals. Examples of successfully cloned mammals include  deer, pigs, and horses.

Serial Nuclear Transplant

One of the major applications of nuclear reprogramming is the production of pluripotent stem cells. The first round of nuclear transplant from differentiated somatic cells into enucleated eggs produces fully viable pluripotent stem cells in scarce proportions. Most embryos from transplanted eggs are not fully reprogrammed; however, grafting these cells into normal embryos increases the proportion of reprogrammed cells. Thus serial nuclear transplantation is a more efficient method of obtaining pluripotent stem cells.

 Core: Cell Biology

Toroids

JoVE 13781

A toroid is a closely wound donut-shaped coil constructed using a single  conducting wire. In general, it is assumed that a toriod consists of  multiple circular loops perpendicular to its axis.

When connected to a supply, the magnetic field generated in the toroid has field lines circular and concentric to its axis. Conventionally, the direction of this magnetic field is expressed using the right-hand rule. If the fingers of the right hand curl in the current direction, the thumb points in the magnetic field direction. The magnetic field inside a toroid varies inversely with the distance from its axis. The magnetic field inside the hollow circle is zero as it does not enclose any current, while outside the toroid, the currents flowing in opposite directions cancel each other out. Hence, the magnetic field is zero.

If a toroid with an inner radius of 10 cm and an outer radius of 15 cm carries a current of 2 A, how many turns does it require to produce a magnetic field of 1 mT 12 cm away from its center? Assume that the turns are equally spaced.

Here, the known quantities are the inner radius, outer radius, current, magnetic field, and distance from the center of the toroid. The number of turns need to be estimated using these known quantities.

The expression for the magnetic field inside a toroid is given as follows:

Equation1

Substituting the values of the known quantities, the number of turns required to produce a magnetic field inside the toroid at a distance of 12 cm from its center is estimated to be 1,500.

 Core: Physics

Uncertainty: Overview

JoVE 14510

In analytical chemistry, we often perform repetitive measurements to detect and minimize inaccuracies caused by both determinate and indeterminate errors. Despite the cares we take, the presence of random errors means that repeated measurements almost never have exactly the same magnitude. The collective difference between these measurements - observed values - and the estimated or expected value is called uncertainty. Uncertainty is conventionally written after the estimated or expected value.

It is important to express the uncertainty with the correct number of significant figures, which is the number of digits required to represent the precise outcome. The magnitude of possible variations from the significant figure in either direction is expressed as addition or subtraction to the significant figure. Uncertainty represented in this way is called absolute uncertainty. The ratio of absolute uncertainty to the magnitude of the estimated or expected value is known as relative uncertainty.

 Core: Analytical Chemistry

Thermodynamics: Activity Coefficient

JoVE 14526

Activity is the measure of the effective concentration of the species in solution. It can be expressed as the product of the molar concentration of the species and its activity coefficient. The activity coefficient is a dimensionless quantity and depends on the total ionic strength of the solution.

The activity coefficient is a measure of the deviation from ideal behavior. When the ionic strength of the solution is minimal, the activity coefficient of an ionic species is close to unity, making activity approximately equal to the molar concentration. Such a solution exhibits behaviors very close to those of an ideal solution. The activity coefficient of an uncharged molecule is approximately unity at all ionic strengths less than 0.1 mol/L.

As the ionic strength of the solution increases, the activity coefficient decreases, making the activity less than the molar concentration. This indicates the deviation of the species from ideal behavior. In solutions with ionic strength higher than 0.1 mol/L, the activity coefficient increases and exceeds unity.

A useful equation for calculating the activity coefficient from ionic strength is the extended Debye–Hückel equation, which relates the coefficient to ionic strength, the effective diameter of the ion, and the ionic charges of the solute. The Debye–Hückel equation relies on the ion size - the effective diameter of a hydrate ion - which can vary widely with charge and the geometry of the ion, thereby introducing uncertainty into the equation.

The equation works well in extremely dilute solutions and solutions with ionic strength less than 0.1 mol/L and extremely dilute solutions, where the uncertainties in ion sizes have negligible effects on the activity coefficients. In these cases, the solution behaves mostly ideally, and the Debye–Hückel equation reduces to a simpler form known as the Debye–Hückel Limiting Law.

 Core: Analytical Chemistry

Titration of a Weak Acid with a Weak Base

JoVE 14542

Weak acids and bases do not undergo dissociation completely, and titrations between these two are rarely studied. When such studies are performed, say, for the titration of a weak acid with a weak base, the titration curve plots the change in pH as a function of the volume of base added. Take the titration of acetic acid with ammonia, for instance. During the titration, these two species form ammonium acetate and water, but the pH change is slow and gradual.

As a result, there is no simple indicator that changes color sharply to help with precisely pinpointing the equivalence point. A relatively simple way to approach this is to mix two indicators – in this case, neutral red and methylene blue – such that the mixture exhibits a sharp color change at the equivalence point – violet-blue to green here.

For a weak acid and weak base that possess identical Ka and Kb values, the pH calculated at the equivalence point is 7. If Ka is greater than Kb, the pH is less than 7, making the endpoint solution acidic. However, if Ka is less than Kb, the calculated pH is greater than 7, resulting in a basic solution at the endpoint.

 Core: Analytical Chemistry

Precipitation Titration: Endpoint Detection Methods

JoVE 14581

In argentometric precipitation titrations, endpoints can be detected visually by the Mohr, Volhard, and Fajans methods. In the Mohr method, adding a soluble chromate indicator gives an initial yellow color to the analyte solution. As the titrant is added, the first excess of silver ions forms a red silver chromate precipitate, marking the endpoint. The solution pH should be maintained at about 8 by adding solid CaCO3.

In the Volhard method, a standard excess of AgNO3 is first added to the chloride ions to precipitate AgCl, which is removed by filtration. Now, the filtrate containing excess Ag+ is back-titrated against thiocyanate in the presence of a ferric ion indicator. A soluble red complex is formed, signaling the endpoint of the titration. The titration must be performed in an acidic medium.

Fajans method employs an anionic adsorption indicator that exhibits different colors when it is in solution and when it is adsorbed. Before the equivalence point, when the chloride ions are in excess, they form the primary adsorbed layer on the AgCl precipitate. The negatively charged precipitate surface repels the anionic indicator, which remains in solution. However, beyond the equivalence point, silver ions are in excess, forming the primary adsorbed layer on the precipitate. The positively charged precipitate surface now attracts the anionic indicator, which is adsorbed. This results in a color change, indicating the endpoint. In this method, the pH of the solution is adjusted based on the indicator used.

 Core: Analytical Chemistry

Other Nuclides: 31P, 19F, 15N NMR

JoVE 14597

Many organic, inorganic, and biological molecules contain spin-half nuclei such as nitrogen-15, fluorine-19, and phosphorus-31. As a result, NMR studies of these nuclei have found extensive applications in chemical and biological research.

While fluorine-19 and phosphorous-31 have high natural abundances (100%) and positive gyromagnetic ratios, nitrogen-15 has a low natural abundance and a negative gyromagnetic ratio. However, nitrogen-15 is still preferred over nitrogen-14 (which has a high natural abundance) because the latter is quadrupolar and produces broad signals. All three nuclei require different standard references and have different chemical shift ranges. For NMR studies of nitrogen-15, nitromethane is used as the standard reference, whereas trichlorofluoromethane and phosphoric acids are used as standard references for fluorine-19 and phosphorus-31 NMR studies, respectively. 

In general, nitrogen-15 chemical shifts are lowest for saturated systems and become increasingly positive with electronegative substitution. However, the lone electron pair on nitrogen can undergo coordination and protonation with the solvent, making the chemical shift solvent-dependent.

While fluorine-19 chemical shifts are also subject to solvent effects, they are predominantly influenced by paramagnetic shielding effects. In addition, fluorine-19 chemical shifts are practically unaffected by ring currents and neighboring group effects. Interestingly, phosphorus-31 chemical shifts are more affected by the coordination number of phosphorus and not strongly influenced by electronegative substituents.

 Core: Analytical Chemistry

Gross Anatomy of Skeletal Muscles

JoVE 14837

The connective tissues play a significant role in arranging the muscle fibers into a hierarchical structure that forms a complete muscle. Consider a muscle like the bicep brachii, commonly called the bicep. This muscle comprises thousands of muscle fibers enclosed by a protective layer of connective tissue called the endomysium. The endomysium is primarily composed of reticular fibers, a type of thin collagen fiber. It allows the exchange of nutrients and waste products at the fiber level, keeping the muscle fibers healthy and functional.

The endomysium is continuous with other layers of connective tissue called the perimysium and epimysium. The perimysium organizes the fibers into groups called fascicles, which are responsible for the synchronized contraction of the bicep muscle during actions such as lifting objects. The organization of fibers into fascicles also helps to distribute the network of nerves and blood vessels more evenly throughout the muscle.

Numerous fascicles are collectively held together by an epimysium, the outermost connective tissue covering the entire muscle. The epimysium also extends from the muscle to form the tendons, which connect the bicep muscle to the bones of the forearm and shoulder. This continuity allows the force generated by muscle contractions to be effectively transferred to the bones, resulting in movement.

These connective tissues, which form the epimysium, perimysium, and endomysium, play a vital role in the structure and function of skeletal muscles. They facilitate muscle contraction and bone movement while supporting and protecting muscle fibers, ensuring their nourishment and overall health.

 Core: Anatomy and Physiology

Muscle Coordination and Action

JoVE 14862

Muscle coordination is a complex and finely tuned process essential for smooth and purposeful movements like flexion, extension, adduction, abduction, and rotation. The human body orchestrates the actions of various muscles working in concert, each with a specific role. Four functional types describe how muscles work together: agonist, antagonist, synergist, and fixator.

Agonists

Agonist muscles, often called prime movers, are the primary muscles responsible for producing a specific movement. For instance, during a bicep curl, the biceps brachii acts as the agonist, contracting to lift the forearm. Agonists are the focal point of any movement, creating the necessary force to change the position of bones and joints.

Antagonists

Antagonist muscles serve as the counterforce to agonists. While the agonist muscle contracts, the antagonist muscle relaxes and lengthens to provide balance and control. This relaxation is crucial for smooth motion and preventing injury. In the bicep curl example, the triceps brachii acts as the antagonist, elongating as the biceps contract. When the movement is reversed and the arm is extended, the roles switch: the triceps become the agonist and the biceps become the antagonist.

Synergists

Synergists aid an agonist by providing additional pull or stabilizing the agonist's origin, their importance varying throughout the movement. Synergists ensure that the motion is smooth and coordinated. For example, during a bicep curl, muscles like the brachialis and brachioradialis work synergistically, providing additional force and ensuring the elbow joint moves correctly.

Fixators

Fixators, or stabilizers, are muscles that hold a part of the body steady, providing a firm base for the agonist to act upon. These muscles might not be directly involved in the movement but are crucial for preventing unwanted motions and maintaining posture. In the bicep curl example, muscles in the shoulder and the upper back act as fixators, keeping the shoulder stable and allowing the biceps to perform the curl efficiently.

 Core: Anatomy and Physiology

Muscles that Move the Leg

JoVE 14880

The movement of the legs is facilitated by numerous muscles located within the anterior, medial, and posterior compartments of the thigh.

Anterior Compartment

The quadriceps femoris, the most visible muscle of the anterior compartment, is integral for leg extension and thigh flexion. It is formed by merging four distinct muscles — the vastus lateralis, vastus medialis, vastus intermedius, and rectus femoris. The quadriceps tendon, a shared tendon of the four quadriceps muscles, is affixed to the patella. This tendon extends beyond the patella, transforming into the patellar ligament that attaches securely to the tibial tuberosity. Additionally, the anterior compartment houses the sartorius muscle, which is involved in knee flexion and lateral hip rotation. It enables sitting cross-legged and is fittingly dubbed the "tailor’s muscle," reflecting its function.

Medial Compartment

The medial compartment of the thigh, also known as the adductor compartment, houses several key muscles, including the gracilis, adductor longus, adductor brevis, adductor magnus, and obturator externus. The gracilis, the most superficial and medial of the thigh muscles, is known for its long, strap-like appearance. It primarily aids in adducting the thigh, flexing the leg, and medially rotating the tibia during knee flexion. The adductor longus, often the most anterior of the adductor muscles, originates from the pubis and extends to the middle third of the linea aspera on the femur. The adductor brevis lies deep and works alongside the adductor longus to adduct and medially rotate the thigh. It has a shorter and more oblique course from the pubis to the femur.

Posterior Compartment

The posterior compartment houses muscles that enable flexion at the knee and extension at the hip, including the three hamstring muscles — the biceps femoris, semitendinosus, and semimembranosus. Hamstring tears or "pulled hamstrings" are common in sports due to their significant involvement in leg movement. The popliteal fossa, a diamond-like area on the backside of the knee, is defined by boundaries set by various tendons. It is flanked on either side by the biceps femoris muscle tendons and medially bordered by the tendons of the semitendinosus and semimembranosus muscles.

Finally, the popliteus is a small, triangular muscle located at the back of the knee. It is known for its crucial role in initiating knee flexion and providing stability to the joint. It originates from the lateral femoral condyle, a bony protrusion on the thigh bone, and then inserts onto the posterior surface of the tibia, just below the knee joint. It contributes to leg movement by causing the medial rotation of the tibia, a key mechanism for unlocking the knee during flexion.

 Core: Anatomy and Physiology

Brainstem

JoVE 14913

The brainstem, located inferior to the brain and superior to the spinal cord, serves as a bridge between the cerebrum and the spinal cord. It plays a vital role in relaying information and controlling critical life functions. It comprises three primary regions: the midbrain, pons, and medulla oblongata.

The Midbrain

The midbrain is located beneath the diencephalon and connects the cerebrum with the lower parts of the brain. The cerebral peduncles are prominent midbrain structures that house the crus cerebri. The crus cerebri contain the corticospinal tract — a bundle of nerve fibers originating from the cerebral cortex. The corticospinal tract carries signals from the brain to the spinal cord, allowing voluntary movements such as walking, reaching, or grabbing objects. Damage to the crus cerebri can result in motor deficits and muscle weakness on the opposite side of the body due to how these nerve fibers cross over in the brainstem.

On its posterior side, the midbrain houses the corpora quadrigemina, consisting of four large nuclei — two superior colliculi involved in visual processing and two inferior colliculi involved in auditory processing. The red nuclei of the midbrain have a rich blood supply and high iron content and are involved in motor coordination, especially the fine motor skills of the limbs. The midbrain also contains the substantia nigra, known for its dopamine-releasing neurons crucial for movement control. When these neurons degenerate, such as with Parkinson’s disease, it results in motor symptoms such as tremors, rigidity, bradykinesia, and postural instability.

The Pons

The pons is a major communication and coordination center between the cerebrum and cerebellum, mainly through the middle cerebellar peduncles. It houses numerous pontine nuclei that transmit signals from the motor cortex to the cerebellum, essential for motor control, coordination, and autonomic regulation. It regulates breathing, contributes to the sleep-wake cycle, and contains cranial nerve nuclei responsible for facial movement and hearing. It also plays a role in coordinating swallowing reflexes and processing taste sensations.

The Medulla Oblongata

The medulla oblongata forms the lower portion of the brainstem and is directly connected to the spinal cord. It features two longitudinal ridges known as the pyramids. They comprise large motor tracts descending from the brain to the spinal cord and play a critical role in controlling voluntary movements. The medulla is a hub for several nuclei that manage both motor and sensory information. It includes the nuclei for the hypoglossal nerve, which controls tongue movements, and vagus nerves, which provide parasympathetic control. The inferior olivary nuclei are involved in motor learning and aid in refining motor skills and coordination. The vestibular nuclei are essential for maintaining balance and gaze stability by processing signals from the inner ear, while the cochlear nuclei process auditory information, enabling hearing and interpretation of sounds from the environment.

 Core: Anatomy and Physiology

Sensory Modalities

JoVE 14931

Sensation typically is the process by which the sensory receptors and sense organs detect stimuli from the internal and external environment and transmit this information to the central nervous system for processing.

General senses refer to the broad category of sensory information detected by receptors in the body and can be further grouped into somatic and visceral senses. Somatic sensations include touch, pressure, temperature, and pain and are essential for navigating our environment and interacting with the world around us. Visceral senses refer to the sensory information from internal body organs, such as the stomach, small and large intestine, and urinary bladder. These sensations, including hunger, thirst, fullness, pain, or discomfort, are essential for regulating internal body functions and maintaining homeostasis.

Special senses are the five sensory modalities dedicated to a specific function and involve specialized sensory organs. These senses are vision, hearing, taste, smell, and balance. Each of these senses relies on specialized receptor cells in specific organs, such as the eyes, ears, tongue, and nose.

 Core: Anatomy and Physiology

Sympathetic Division of the ANS

JoVE 14948

The sympathetic division of the autonomic nervous system (ANS) plays a crucial role in preparing the body for stress, physical activity, and increased energy demands. This division activates the "fight-or-flight" response, enabling individuals to respond effectively to challenging situations.

Originating in the thoracic and lumbar spinal cord segments, the preganglionic fibers of the sympathetic division exit the spinal cord through the white ramus communicans. They then enter the sympathetic trunk ganglion, where they have three possible locations for synapse with postganglionic fibers: within the sympathetic chain ganglia, at the collateral ganglia, or in the adrenal medulla.

Sympathetic Chain Ganglia

The sympathetic chain ganglia, also known as the paravertebral ganglia, are a series of interconnected ganglia positioned bilaterally along the spinal cord. These ganglia form a chain-like structure that runs parallel to the vertebral column. The sympathetic chain ganglia play a crucial role in regulating various physiological functions throughout the body.

Collateral Ganglia

The collateral ganglia, also known as the prevertebral or preaortic ganglia, are positioned anteriorly to the spinal cord. Once within the collateral ganglia, the preganglionic fibers synapse with the postganglionic fibers. These postganglionic fibers extend from the ganglia and innervate specific target organs in the abdominal and pelvic regions. Organs innervated by the collateral ganglia include the heart, lungs, liver, stomach, intestines, and kidneys. The ganglia are involved in the control of functions such as cardiac activity, bronchial dilation, gastrointestinal motility, and urinary system regulation.

Adrenal Medulla

The adrenal medulla, found within the adrenal gland, plays a unique role in the sympathetic division. When stimulated, it releases neurotransmitters such as adrenaline and noradrenaline into the bloodstream. These neurotransmitters act as hormones and regulate the activity of distant cells, further enhancing the body's response to stress and physical demands.

 Core: Anatomy and Physiology

Taste Buds and Receptors

JoVE 14967

Gustation, or the sense of taste, is intrinsically linked to the anatomical structures located on the tongue. This organ's surface, along with the entirety of the oral cavity, is adorned with stratified squamous epithelium. Evident on the tongue are elevated structures known as papillae (singular = papilla), which house the mechanisms for the transduction of gustatory stimuli. Four distinct types of papillae exist, each identified by their unique morphological attributes: the circumvallate, foliate, filiform, and fungiform papillae. Nestled within the papillae architecture are taste buds, which harbor specific gustatory receptor cells dedicated to the transduction of taste stimuli. These cells exhibit sensitivity to the chemical composition of ingested food, releasing neurotransmitters commensurate with the chemical concentration of the food. The released neurotransmitters can stimulate sensory neurons in the cranial nerves: facial, glossopharyngeal, and vagus.

The human gustatory system, responsible for perceiving taste, is an intricate matrix of processes that combine to discern food flavors. On the microscopic level, gustatory cells are activated by chemical constituents in our food, triggering the release of neurotransmitters. These chemicals relay signals to the brain, where taste perception, such as 'sweet,' is formed. Despite its rapid execution, this system's complexity is evident in its various stages: chemical detection by gustatory receptors, accurate and swift signal transmission via nerves, and efficient information processing by the brain. Furthermore, taste perception is not a solitary phenomenon but is influenced by other sensory inputs like olfactory cues and tactile sensations. For instance, sweet-smelling food might be perceived as bitter if the texture is undesirable. This exemplifies the gustatory system's remarkable intricacy, hinting at integrating multiple sensory modalities to generate a precise flavor profile. Besides contributing to the hedonic enjoyment of food, taste perception plays a crucial role in assessing food safety. A comprehensive understanding of the biology and sensory perception surrounding taste can provide insights into our body's food processing mechanisms and flavor experiences, aiding in making informed, health-conscious dietary selections and enhancing our appreciation of taste subtleties.

Commonly Known Gustatory Dysfunction

1) Ageusia refers to the complete loss of taste, while hypogeusia is a reduced ability to taste sweet, sour, bitter, salty, and umami flavors. Conversely, dysgeusia is a condition characterized by a distorted sense of taste.

Common symptoms of these disorders include difficulty identifying specific tastes or experiencing an unpleasant or strange taste in the mouth. They may also cause changes in appetite and dietary habits, potentially leading to weight gain or loss.

2) Phantogeusia is a less common but equally impactful gustatory disorder. It involves tasting something that is not present, often described as a persistent, phantom flavor. This can be particularly distressing and disruptive to daily life, making eating and drinking unpleasant experiences.

Causes of Gustatory Dysfunction

The causes of gustatory disorders are diverse. They range from aging, which naturally diminishes taste sensitivity, to more specific factors, such as medication side effects, infections, and head injuries. Other potential causes include radiation therapy, certain neurological conditions like Parkinson's or Alzheimer's disease, and nutritional deficiencies.

Treatment Options

Treatment for gustatory disorders largely depends on the underlying cause. A doctor may adjust the dosage or switch to a different drug if medication is the culprit. Infections may be treated with antibiotics, while nutritional deficiencies can be addressed through diet modifications or supplements.

Surgical interventions may be considered in some cases, particularly if a nerve injury or tumor is causing the taste disorder. However, these are usually last-resort options due to their invasive nature and potential risks.

Pharmacological interventions are another treatment avenue. These may involve drugs that aim to enhance taste sensation or manage symptoms. For instance, zinc supplements have shown promise in improving taste function in some individuals.

The Future of Gustatory Disorder Treatment

Research into gustatory disorders is ongoing, with studies exploring new treatment options and seeking to understand these conditions' underlying mechanisms better. While progress has been made, much remains to be discovered. By continuing to advance our knowledge in this area, we hope to improve the lives of those living with these challenging disorders.

In conclusion, gustatory disorders are complex conditions that can significantly impact an individual's ability to enjoy food and maintain a healthy diet. Understanding their causes, symptoms, and treatments is essential for providing adequate care and improving patients' quality of life.

 Core: Anatomy and Physiology

Functions of Thyroid Hormones

JoVE 14983

The thyroid hormone (TH) plays a pivotal role in the intricate orchestration of physiological processes, exerting profound effects on development, metabolism, and homeostasis throughout different life stages.

TH is indispensable for the normal development and maturation of the skeletal, muscular, and nervous systems during fetal and childhood growth. It facilitates bone mineral turnover and regulates protein synthesis in developing tissues, contributing significantly to overall growth and maturation.

TH continues its influence upon reaching adulthood, particularly on the reproductive and circulatory systems. Primarily in females, it regulates gonadal hormone secretion, impacting the reproductive system, and stimulates endocrine tissue, contributing to hormonal balance. In the circulatory system, TH promotes normal heart function, facilitates red blood cell synthesis, and enhances tissue oxygenation.

TH binds to various receptors within the cell, including cytoplasmic, nuclear, and mitochondrial receptors. Cytoplasmic receptors store TH and release it when intracellular TH levels decrease, maintaining cellular homeostasis. TH activates glycolysis and ATP production genes in the mitochondria, increasing energy synthesis. Simultaneously, TH indirectly enhances ATP production in the nucleus by controlling the synthesis rate of enzymes involved in the cellular metabolic rate.

The culmination of these actions is known as the calorigenic effect, reflecting an elevated cellular metabolic rate induced by thyroid hormone. This multifaceted role underscores the critical contribution of TH to fundamental physiological processes, ensuring the harmonious functioning of various organ systems and the maintenance of overall cellular homeostasis.

 Core: Anatomy and Physiology

Independent and Dependent Sources

JoVE 15068

In electrical circuits, sources play a crucial role in providing power for the operation of the circuit. These sources can be broadly categorized into two types: independent and dependent.

Independent voltage or current sources supply a fixed amount of voltage or current, respectively, which is unaffected by other elements within the circuit. These are represented using specific symbols. Independent voltage sources are symbolized with polarities (+ and -), indicating the direction of the potential difference, while independent current sources use arrows to denote the direction of the current flow.

An interesting aspect of independent sources is that a short circuit can be considered as a zero-voltage independent voltage source because it allows current to flow freely with no voltage drop. Conversely, an open circuit can be thought of as a zero-current independent current source, given that it prevents any current flow regardless of the applied voltage.

On the other hand, dependent sources are voltage or current generators whose values are controlled by other elements within the same circuit. They are depicted by a diamond shape in circuit diagrams.

These dependent sources are further divided into four categories based on their controlling and controlled parts: voltage-controlled voltage sources (VCVS), current-controlled voltage sources (CCVS), voltage-controlled current sources (VCCS), and current-controlled current sources (CCCS).

For instance, in a circuit where the voltage output of a dependent source is controlled by the current flowing through a resistor, if a short circuit is added in series with this resistor, the controlling current becomes equivalent to the current in the short circuit. The ratio of the source voltage to the controlling current represents the gain of the dependent source.

Similarly, the controlling voltage can be considered as open-circuit voltage. This means that the voltage or current of the dependent sources is proportional to the controlling voltage or current.

 Core: Electrical Engineering

Capacitors

JoVE 15084

Capacitors play a crucial role in car radios, where they filter and store frequencies to ensure clear signal reception. Essentially serving as energy storage devices, capacitors store energy within their electric field and are composed of two parallel conducting plates separated by a dielectric.

When a voltage source is connected to a capacitor, positive and negative charges accumulate on the opposite plates. This accumulation generates a potential difference that equals the product of the electric field and the distance between the plates. This process continues until the potential difference reaches the source voltage.

The electric field within a capacitor is proportional to the charge stored on the plates and the area of the plates. The stored charge is directly proportional to the applied voltage. The proportionality constant in this relationship is known as capacitance, which indicates the amount of charge needed to create a given potential difference. The unit used to measure capacitance is the farad.

For parallel plate capacitors, the capacitance is directly proportional to the plate area and the permittivity of the dielectric, which is the material's ability to store electrical energy. Meanwhile, the capacitance is inversely proportional to the distance between the plates; the closer the plates, the higher the capacitance.

By differentiating the charge-voltage equation of a capacitor with respect to time, one can derive the current flowing through the capacitor. This is particularly important when considering how capacitors behave in circuit designs.

Furthermore, when a charged capacitor is connected to a load, it discharges. During this discharge process, electrons flow in the reverse direction until the potential difference across the plates reaches zero. This ability to store and release electrical energy makes capacitors invaluable components in many electronic devices. By understanding how they work, engineers can design more effective and efficient electronic systems.

 Core: Electrical Engineering

Parallel RLC Circuits

JoVE 15101

Street lamps equipped with RLC surge protectors are an excellent example of applying circuit analysis in practical scenarios. These surge protectors safeguard the lamp's components against sudden voltage spikes.

A simplified parallel RLC circuit model with a DC input source generating a step response is employed in this context. When the switch is turned on, Kirchhoff's current law is applied, leading to a second-order differential equation.

Equation1

Interestingly, this equation's solution comprises both transient and steady-state responses. The transient response gradually diminishes over time, exhibiting similarities to the source-free series RLC circuit solutions under different damping conditions. If the damping factor surpasses the resonant frequency, the response becomes overdamped.

Equation2

When these two factors match, the response is critically damped.

Equation3

The response turns underdamped if the damping factor is less than the resonant frequency.

Equation4

The steady-state response corresponds to the final inductor current, aligning with the source current. Determining the constants involved relies on the initial conditions of the circuit.

Notably, only the transient response remains when the input source current is eliminated. Parallel RLC circuits find extensive applications, particularly in communications networks and filter designs. Understanding their behavior under different damping scenarios contributes to adequate surge protection and circuit design in practical settings.

 Core: Electrical Engineering

Design Example: Strain Gauge Bridge or Wheatstone Bridge

JoVE 15171

The utilization of strain gauges as transducers for converting mechanical strain into electrical signals is a common practice in various engineering applications. These strain gauges are frequently integrated into Wheatstone bridge circuits to accurately measure parameters such as force or pressure. Within this context, each element within the circuit exhibits a resistance that undergoes subtle variations when subjected to mechanical strain. The primary objective is to convert minuscule variations in small voltage output into a more discernible voltage output, which can be conveniently read using a voltmeter.

In the engineer's capacity, the responsibilities include creating a strain gauge design and determining the necessary amplification to reflect changes in resistance accurately. This relies upon the application of Thevenin's theorem, which establishes the connection between the output voltage of the bridge circuit and variations in resistance.

The procedure starts by calculating the Thevenin voltage, which requires the initial determination of currents within both the upper and lower branches of the circuit. As a result, these computed current values are substituted into Ohm's law to derive the Thevenin voltage. Simultaneously, Thevenin resistance is determined by removing the voltage source in a simplified configuration.

The Thevenin equivalent circuit finds the bridge's output voltage through the voltage division rule. The critical outcome of this analysis enables the precise determination of the amplifier gain required for designing the circuit to operate effectively within the specified operational range.

 Core: Electrical Engineering

Neuronal Transfection Methods

JoVE 5215

Transfection - the process of transferring genetic material into cells - is a powerful tool for the rapid and efficient manipulation of gene expression in cells. Because this method can be used to silence the expression of specific proteins or to drive the expression of foreign or modified proteins, transfection is an extremely useful tool in the study of the cellular and molecular processes that govern neuron function. However, mature neurons have a number of properties that make them difficult to transfect, so specialized techniques are required for the genetic manipulation of this cell type.

This video reviews the principles and rationale behind transfecting neurons. Three common strategies for neuronal transfection are discussed, including nucleofection, gene-gun, and viral transduction. In addition to describing how each of these techniques overcomes the challenges associated with transfecting neurons, the presentation includes a description of how all three methods are performed. Finally, several applications of neuronal transfection are introduced, such as the expression of fluorescent tubulin proteins to visualize neuron morphology, and selective gene silencing to generate a cell culture model of Parkinson’s disease.

 Neuroscience

An Introduction to Aging and Regeneration

JoVE 5337

Tissues are maintained through a balance of cellular aging and regeneration. Aging refers to the gradual loss of cellular function, and regeneration is the repair of damaged tissue generally mediated by preexisting adult or somatic stem cells. Scientists are interested in understanding the biological mechanisms behind these two complex processes. By doing so, researchers may be able to use somatic stem cells to treat degenerative diseases and develop therapies that could delay the effects of aging.

In this video, we provide a brief history of the field of aging and regeneration, touching upon observations made in ancient Greece, as well as modern-day experiments. Some of the questions being asked in this field, and the prominent methods being used by biologists to answer them, are then explored. Finally, we look at a few specific experiments being conducted in today\'s aging and regeneration research laboratories.

 Developmental Biology

An Introduction to Modeling Behavioral Disorders and Stress

JoVE 5428

Recently, it has been discovered that behavioral conditions such as, depression, anxiety and stress-response have a neurological basis. Understanding the biological underpinnings of these conditions may help scientists in developing more effective therapies to treat these disorders. Typically, rodent models are used in this field and behavioral scientists create these models using interventions like drug administration or surgery. It is important to understand how to create and evaluate rodent models of behavioral disorders as they play an important role in discovery of new treatments for clinical applications.

Here, JoVE science education video first reviews the 'classical' criteria used to evaluate rodent models of stress and behavioral disorders. This is followed by some of the important questions that scientists are trying to answer using these models as tools. We'll also go over several rodent behavioral tests currently being used in this field and discuss applications of these paradigms.

 Behavioral Science

An Overview of Genetics and Disease

JoVE 5543

Many human diseases are associated with mutations or variations in genetic sequences. Some of these genetic variants are heritable, passed down from generation to generation, while others arise sporadically during an organism’s life and cause diseases such as cancer. Researchers are trying to identify and characterize these genetic alterations in the hopes of improving diagnosis and therapeutic options for patients.

In this video, we will examine the history of genetic disease research, and explore key questions asked by medical geneticists. Various tools used to identify the genetic basis of diseases are then discussed, including genotyping techniques and genome-wide association studies (GWAS). Finally, several current examples of medical genetics research are presented.

 Genetics

Introduction to Mass Spectrometry

JoVE 5634

Source: Laboratory of Dr. Khuloud Al-Jamal - King's College London

Mass spectrometry is an analytical chemistry technique that enables the identification of unknown compounds within a sample, the quantification of known materials, the determination of the structure, and chemical properties of different molecules.

A mass spectrometer is composed of an ionization source, an analyzer, and a detector. The process involves the ionization of chemical compounds to generate ions. When using inductively coupled plasma (ICP), samples containing elements of interest are introduced into argon plasma as aerosol droplets. The plasma dries the aerosol, dissociates the molecules, and then removes an electron from the components to be detected by the mass spectrometer. Other ionization methods such as electrospray ionization (ESI) and matrix assisted laser desorption ionization (MALDI) are used to analyze biological samples. Following the ionization procedure, ions are separated in the mass spectrometer according to their mass-to-charge ratio (m/z), and the relative abundance of each ion type is measured. Finally, the detector commonly consists in an electron multiplier where the collision of ions with a charged anode leads to a cascade of increasing number of electrons, which can be detected by an electrical circuit connected to a computer.

In this video, the procedure of ICP-MS analysis will be described by the detection of 56Fe as an example.

 Analytical Chemistry

Detecting Reactive Oxygen Species

JoVE 5654

Reactive oxygen species are chemically active, oxygen-derived molecules capable of oxidizing other molecules. Because of their reactive nature, there are many deleterious effects associated with unchecked ROS production, including structural damage to DNA and other biological molecules. However, ROS can also be mediators of physiological signaling. There is accumulating evidence that ROS play significant roles in everything from activation of transcription factors to the mediation of inflammatory toxicity that kills foreign pathogens and defend the body.

In this video we will delve into the associations between ROS, metabolism and disease. After establishing their significance, we will discuss the principles and a protocol of a commonly used methodology for measuring ROS levels in cells: the use of non-fluorescent probes that become fluorescent upon oxidation. Lastly, we will review some current applications of this technique in cell biology research.

 Cell Biology

Surface Plasmon Resonance (SPR)

JoVE 5697

Surface plasmon resonance (SPR) is the underlying optical phenomenon behind label-free biosensors to evaluate the molecular affinity, kinetics, specificity, and concentration of biomolecules. In SPR, biomolecular interactions occur on a biosensor made of a thin layer of metal on a prism. Real-time interactions of biomolecules can be monitored by measuring the changes of light reflected off the underside of the metal.

This video describes the basic concepts of SPR and how it is used to analyze and visualize biomolecular interactions. This is followed by a sample preparation and experimental protocol for investigating binding rates using SPR. In the applications section, SPR imaging, localized SPR, and quantum dot enhanced SPR are explored.

Surface plasmon resonance, or SPR, is the underlying phenomenon behind certain label-free biosensors for evaluating binding and adsorption interactions of biomolecules. Binding assays that require labeling, such as ELISA, can be a time-consuming process, and may alter the functionality of the analyte. In SPR, biomolecular interactions occur on a special sensor made of a thin layer of metal on one face of a prism. By monitoring the changes in light reflected off of the underside of the metal, SPR instruments visualize these interactions in real-time without the use of labels. This video will introduce the principles of SPR, a general procedure for SPR imaging, and some applications of in biochemistry.

An SPR sensor is usually made of a thin layer of a noble metal atop the face of a prism. To take readings from the sensor, light is reflected off of the prism-metal interface into a photodetector. The reflected light will have a high intensity except at a certain angle, related to the electronic properties of the metal surface, known as the "surface plasmon resonance angle".

As molecules bind to the surface, the electronic properties of the metal change, which in turn adjusts the angle. As new proteins attach, forming complexes, the angle will shift further. By measuring relative changes in the SPR angle, interactions like these can be monitored in real-time.

Another technique called localized, or "L"SPR, uses metal nanoparticles as the sensor surface. The properties that affect the SPR angle are highly localized to each nanoparticle, which improves sensitivity and signal resolution.

When investigating binding interactions with standard SPR, the sensor is generally mounted in a platform that becomes the floor of a flow cell in the instrument. The biomolecules of interest are carried through the flow cell by buffer solution. The sensor surface is often first coated with a substrate that has a high affinity for the metal. This ensures that a significant amount of ligand, which in turn binds to the analyte of interest, will be immobilized onto the sensor and reduces the likelihood that the ligand will dissociate during the procedure.

Once the ligand is immobilized on the sensor, the analyte is flowed over the sensor in buffer. By monitoring the change in the SPR angle over time as the analyte binds to the ligand, the binding rate and other kinetic information can be calculated.

The reflectance data can also be used for SPR imaging, or SPRi, by directing the reflected light to a CCD detector. This produces a high-contrast, high-resolution image of the entire sensor surface. Using SPR and the related techniques, questions can be answered about molecular affinity, kinetics, specificity, and concentration.

Now that you understand what is being measured in an SPR experiment, let's look at a procedure for investigating binding rates.

Before beginning the procedure, the running and sample buffers must be prepared. The running buffer is used to deposit the ligand onto the sensor, and the sample buffer is used to deposit the analyte. The sensor chip is carefully cleaned and loaded into a sheath. Then, the device is placed into the instrument, where it becomes the bottom of the flow cell. The instrument software is set up for the experiment and subsequent analysis. If necessary, the sensor surface is primed with a substrate to capture the ligand. The ligand is flowed over the sensor surface in the running buffer, where it is captured by the substrate on the sensor surface.

Then, the analyte in the sample buffer is run through the flow cell, where it selectively binds to the immobilized ligand. The change in reflectance is plotted and compared against controls to determine rate constants and other reaction kinetics data for the investigated reaction.

Now that you understand how an SPR experiment is performed, let's look at a few other applications of SPR in biochemistry.

Here, SPR imaging was used to evaluate proteins with an array of eleven receptors on a sensor. 3D graphs of reflectivity versus time and receptor concentration were prepared from the reflectivity data for each protein. These "profiles" are characteristic to each protein, and thus could subsequently be used for protein identification.

In this experiment, cell secretions were studied using a custom-made LSPR sensor. The sensor was also compatible with SPRi and fluorescence microscopy. Upon depositing the cell on the sensor, the interaction of cell secretions with the nanoparticle array could be measured with high spatial resolution.

Here, the use of quantum dots, nanoscale semiconductors, as an SPR signal enhancement agent mixed with the analyte was investigated. This enhanced "nano-SPRi" method was compared to assays by standard SPRi and the ELISA method. The nano-SPRi method significantly improved the sensitivity and limit of detection while still being less time-consuming than the ELISA method.

You've just watched JoVE's video on surface plasmon resonance. This phenomenon is used to monitor and image biomolecular interactions without the use of labels. This video introduced the principles of SPR, a typical protocol for performing an SPR experiment, and a few applications of SPR in biochemistry.

Thanks for watching!

 Biochemistry

Electrochemical Biosensing

JoVE 5796

Electrochemical biosensors detect the binding of a target molecule by sensing an oxidation-reduction event. These sensors paved the way for modern biosensing after the invention of the glucose biosensor. This video will introduce electrochemical biosensing, show the workings of the glucose biosensor, and discuss how electrochemical biosensors are used in cancer detection.

 Bioengineering

Conformations of Cyclohexane

JoVE 11712

Cyclohexane does not exist in a planar form due to the high angle and torsional strain it would experience in the planar structure. Instead, it adopts non-planar chair and boat conformations.

The chair form is the most stable and derives its name from its resemblance to the “easy chair.” In the chair conformation, two carbon atoms are arranged out-of-plane — one above and one below, minimizing the torsional strain. In the chair form, the bond angle is very close to the ideal tetrahedral value, and hence, this form does not show any angle strain. The boat conformation, which derives its name from the resemblance to a boat, has 30 kJ/mol more strain energy than the stable chair form. Although the boat form is devoid of the angle strain, it has considerable torsional strain due to the upward-facing methylene groups. Additionally, due to the proximity of the hydrogen atoms on these groups, strong van der Waals repulsions, known as “flagpole interactions,” further destabilizes the boat form. Hence, the boat form slightly twists one of its C-C bonds to create the twist-boat conformation. Due to the twist, the flagpole hydrogens are now placed slightly apart, reducing the overall strain by 7 kJ/mol.

 Core: Organic Chemistry

Structure and Bonding of Alkenes

JoVE 11764

Olefins, which are unsaturated hydrocarbons containing one or more carbon–carbon double bonds, are broadly divided into alkenes and cycloalkenes. The general chemical formula of an alkene is CnH2n.

Doubly bonded carbons are sp2 hybridized and have a trigonal planar geometry. The double bond is composed of a σ bond formed by the overlap of hybrid orbitals and a π bond produced by the lateral overlap of unhybridized 2p orbitals on both the carbons. Each carbon atom is bonded to two hydrogen atoms through sp2s orbital overlap. As the unhybridized p electrons have higher energy than the electrons in the hybrid orbitals, the π bond electrons generally have higher energy than the σ bond electrons. Therefore, electrophiles react preferentially with the π bond electrons of alkenes.

Ethylene, propylene, and butylene exist as colorless gases. Alkenes containing 5 to 14 carbon atoms are liquids, and those containing 15 or more carbon atoms are solids. Alkenes, being effectively nonpolar, are insoluble in water but soluble in nonpolar solvents.

The boiling points of alkenes increase with the molecular mass as the intermolecular forces become stronger with the increasing size of the molecules.

Table 1: Physical Properties of Alkenes and Cycloalkenes

Name  Structural Formula mp (°C)  bp (°C)
Ethylene CH2=CH2 −169 −104
Propylene CH3CH=CH2 −185 −47
1-Butene CH3CH2CH=CH2 −185 −6
1-Pentene CH3CH2CH2CH=CH2 −165 30
Cyclopentene C5H8 −135 44
Cyclohexene C6H10 −104 83

Alkenes occur abundantly in nature. For example, ethylene, the simplest alkene, is found in nature as a plant hormone that affects the ripening of fruits. Lycopene and carotenes are the polyenes responsible for the red, orange colors of fruits and vegetables, like tomatoes and carrots. Also, alkenes are the structural frames of various plant essential oils and insect pheromones.

 Core: Organic Chemistry

Oxidation of Alkenes: Anti Dihydroxylation with Peroxy Acids

JoVE 11783

Diols are compounds with two hydroxyl groups. In addition to syn dihydroxylation, diols can also be synthesized through the process of anti dihydroxylation. The process involves treating an alkene with a peroxycarboxylic acid to form an epoxide. Epoxides are highly strained three-membered rings with oxygen and two carbons occupying the corners of an equilateral triangle. This step is followed by ring-opening of the epoxide in the presence of an aqueous acid to give a trans diol. Peroxycarboxylic acids are strong oxidizing agents and analogous to carboxylic acids. However, they have an extra oxygen atom between the carbonyl group and the hydrogen atom. Commonly used organic peracids include meta-chloroperoxybenzoic acid and peroxyacetic acid.

The mechanism begins with a concerted nucleophilic attack by the alkene π bond on the electrophilic oxygen of the peroxy acid, breaking the oxygen–oxygen bond and forming a new carbon–oxygen double bond, leading to a cyclic transition state.

Figure1

The second step of the reaction involves an acid-catalyzed ring-opening of the epoxide to finally form a trans diol.

Figure2

Overall, the regiochemistry of the reaction is governed by a combination of steric and electronic factors. In epoxides with a primary and secondary carbon, steric factors dominate, favoring an attack at the less substituted carbon. With a tertiary carbon, electronic effects dominate, favoring attack at the more substituted carbon.

 Core: Organic Chemistry

Reduction of Alkynes to cis-Alkenes: Catalytic Hydrogenation

JoVE 11841

Introduction

Like alkenes, alkynes can be reduced to alkanes in the presence of transition metal catalysts such as Pt, Pd, or Ni. The reaction involves two sequential syn additions of hydrogen via a cis-alkene intermediate.

Figure1

Thermodynamic Stability

Catalytic hydrogenation reactions help evaluate the relative thermodynamic stability of hydrocarbons. For example, the heat of hydrogenation of acetylene is −176 kJ/mol, and that of ethylene is −137 kJ/mol. The higher exothermicity associated with the addition of hydrogen to acetylene suggests that it is thermodynamically less stable than ethylene.

Figure2

Figure3

Modified Catalyst

Catalytic reduction of alkynes can be stopped at the cis-alkene stage using a modified or poisoned catalyst such as Lindlar or P-2 catalyst. The Lindlar catalyst consists of palladium metal deposited on calcium carbonate and modified using lead acetate and quinoline; the P-2 catalyst is a nickel-boride complex.

Figure4

A modified catalyst lowers the activation energy for the reduction of the first π bond; however, it is not powerful enough to catalyze the reduction of the second π bond. For example, hydrogenation of 2-pentyne over Lindlar catalyst gives cis-2-pentene.

Figure5

Hydroboration-Protonolysis

Hydroboration-protonolysis is a non-catalytic method for the conversion of internal alkynes into cis-alkenes. The reaction involves treating an internal alkyne with borane to form a trialkenylborane intermediate followed by treatment with acetic acid to yield the desired cis-alkene.

Figure6

 Core: Organic Chemistry

Preparation of Alcohols via Addition Reactions

JoVE 11923

Overview

The acid-catalyzed addition of water to the double bond of alkenes is a large-scale industrial method used to synthesize low-molecular-weight alcohols. An acidic atmosphere is required to allow the hydrogen in the water molecule to act as an electrophile and attack the double bond in an alkene. The addition of a proton to the double bond creates a carbocation intermediate. The proton preferentially bonds to the less substituted end of the double bond to create a more stable carbocation intermediate. Subsequently, water attacks the carbocation to yield a  protonated alcohol.

Here, the reaction is regioselective and prefers Markovnikov's addition of water to the alkene. The proton is added to the less substituted end of the double bond, and the -OH group is added to the more substituted carbon on the other end. The reaction is not stereoselective since the water can approach the planar carbocation from both sides, which leads to a racemic mixture of alcohols. Water addition can happen from the same side (syn addition) or the opposite side (anti addition) with respect to the direction from which the proton is added to the double bond. In alkenes with a more substituted carbon adjacent to the double bond, the positive charge in the carbocation intermediate may rearrange (1,2 hydride shift or 1,2 methyl shift) to gain stability. In this situation, a mixture of constitutional isomers of alcohols is formed as products.

An alternative synthetic route for preparing alcohols from alkenes is via oxymercuration-demercuration reaction. In the oxymercuration reaction, the electrophilic mercury in mercury acetate is added to the double bond and forms a three-membered ring structure called mercurinium ion. Water then attacks the mercurinium ion to form organomercurial alcohol. The subsequent step is demercuration, where sodium borohydride replaces the mercury species with hydrogen.

The hydroxyl group bound to the more substituted carbon in the mercurinium ion guides the hydrogen towards the less substituted carbon, eventually resulting in a Markovnikov's product. Also, attack by water is anti with respect to mercury acetate addition. However, the reaction that replaces mercury species with hydrogen is not stereoselective, which results in the formation of both enantiomers of the alcohol as the final product. Due to the absence of a conventional carbocation and its rearrangement, oxymercuration-demercuration has better yields than direct acid-catalyzed hydration.

Anti-Markovnikov's addition of water to an alkene can be performed by the hydroboration-oxidation method. Hydroboration is the single-step addition of a B-H bond from a borane to the two carbons in the double bond via a four-membered transition state. Electrophilic boron attracts the π orbital electrons from the least hindered carbon in the double bond. The resulting partial positive charge is stabilized on the other more substituted carbon atom. Steric factors also favor the addition of boron to the least hindered carbon and hydrogen to the other. The alkyl borane intermediate further reacts with two more alkene molecules to yield a trialkyl borane intermediate. Oxidation of trialkyl borane with hydrogen peroxide in sodium hydroxide replaces the boron with a hydroxyl group. The simultaneous addition of boron and hydrogen from borane ensures a syn addition over the double bond. However, the syn addition can happen from either face of the alkene.

Due to the difference in regiochemistry and reaction mechanisms, these three different synthetic pathways result in different alcohols from the same alkene substrate, as shown in

Figure1

Figure 1. Addition of water over 3,3-dimethyl-1butene under acid-catalysis (top),  oxymercuration-demercuration (middle) and hydroboration-oxidation conditions (bottom). 

 Core: Organic Chemistry

Phagocytosis of Apoptotic Cells

JoVE 12430

Cells undergoing apoptosis form apoptotic bodies that must be removed immediately to prevent inflammation, autoimmune diseases, and necrosis. Phagocytosis is carried out by professional phagocytes such as macrophages or  immature dendritic cells. Non-professional phagocytes such as  epithelial cells and fibroblasts also take part in this process; however, they are not as effective as professional phagocytes. 

Normal cells contain receptors that prevent them from being recognized by phagocytes. For example, red blood cells display CD47 receptors that block phagocytosis. Another marker, CD31, also known as platelet endothelial cell adhesion molecule-1, is expressed by normal cells, which prevents phagocyte attachment. During apoptosis, CD31 is not expressed, leading to phagocyte recognition and phagocytosis.

Apoptotic cells release 'find-me' signals such as  lysophosphatidylcholine, sphingosine-1-phosphate, fractalkine, thrombospondin-1, ATP, and UTP to help phagocyte recognition. Once the phagocytes reach the location, the apoptotic cells display 'eat-me' signals, such as apoptotic cell-associated molecular patterns (ACAMP) and phosphatidylserine. When phagocytes attach to apoptotic cells, they release anti-inflammatory cytokines such as IL-10 and TGF-ꞵ, creating a non-inflammatory microenvironment. This prevents damage to neighboring cells. Once the apoptotic cells are internalized, they are broken down into nucleotides, amino acids, and sterols which are further recycled.

 Core: Cell Biology

Osteoclasts in Bone Remodeling

JoVE 12520

Osteoclasts are cells responsible for bone resorption and remodeling. They originate from hematopoietic progenitor cells present in the bone marrow. Numerous progenitor cells fuse to form multinucleated cells, each with 10-20 nuclei. A single osteoclast has a diameter of 150 to 200 µM. These cells have ruffled borders that break down the underlying bone tissue and release minerals such as calcium into the blood in bone resorption. Osteoclasts cling to bones with their ruffled edges during bone resorption and secrete several enzymes, including acid phosphatase. The acid phosphatase mineralizes the bone by degrading organic collagen and releasing calcium and phosphorus.

Bone remodeling is a skeletal change that occurs regularly in coordination with bone formation. This balances the breakdown and formation of new bone. However, after around 40 years of age, bone resorption occurs more frequently than formation, resulting in a reduction in bone density. This results in osteoporosis, which makes the bones weaker and brittle, increasing the risk of fractures.

Several hormones and proteins regulate the process of bone resorption. The Receptor activator of nuclear-factor kappa or RANK binds to its ligand and stimulates bone resorption. Calcitonin, a hormone released by the thyroid gland, reduces circulating calcium in the blood and inhibits bone resorption, thereby promoting bone formation. Similarly, parathyroid hormone or PTH increases the calcium level in the blood. Additionally, it also increases the activity of RANKL, promoting bone resorption. Growth hormone stimulates the activity of both osteoblasts and osteoclasts, thereby enabling both bone resorption and formation simultaneously.

The hormone estrogen negatively regulates bone resorption. A deficiency in estrogen increases bone resorption and bone remodeling. Calcitonin inhibits the resorption process by binding to calcitonin receptors on osteoclasts. Thus, calcitonin plays a role in calcium homeostasis and is used to treat osteoporosis.

 Core: Cell Biology

Plant Tissues

JoVE 13364

Plants are multicellular eukaryotes with tissue systems made of various cell types that carry out specific functions. Different tissues work together to perform a unique function and form an organ. Organs working together form organ systems. Vascular plants have two distinct organ systems: a shoot system and a root system. The shoot system consists of two portions: the vegetative (non-reproductive) parts of the plant, such as the leaves and the stems, and the reproductive parts of the plant, which include flowers and fruits. The shoot system generally grows above ground, absorbing the light needed for photosynthesis. The root system, which supports the plants and absorbs water and minerals, is usually underground.

Plant tissue systems comprise meristematic and permanent (or non-meristematic) tissue. Cells of the meristematic tissue are found in meristems, which are plant regions of continuous cell division and growth. Meristematic tissue cells are either undifferentiated or incompletely differentiated and continue to divide and contribute to the plant's growth. In contrast, permanent tissue consists of plant cells that are no longer actively dividing.

Meristematic tissues consist of three types, based on their location in the plant. Apical meristems contain meristematic tissue located at the tips of stems and roots, which enable a plant to extend in length. Lateral meristems facilitate growth in thickness or girth in a maturing plant. Intercalary meristems occur only in monocots, at the bases of leaf blades, and at nodes (the areas where leaves attach to a stem). This tissue enables the monocot leaf blade to increase in length from the leaf base; for example, it allows lawn grass leaves to elongate even after repeated mowing.

Secondary tissues are either simple (composed of similar cell types) or complex (consisting of different cell types). Dermal tissue, for example, is a superficial tissue that covers the outer surface of the plant and controls gas exchange. Vascular tissue is an example of a complex tissue made of two specialized conducting tissues: xylem and phloem. Xylem tissue transports water and nutrients from the roots to different plant parts. It includes three cell types: vessel elements, tracheids (both of which conduct water), and xylem parenchyma. Phloem tissue transports organic compounds from photosynthesis to other plant parts and consists of four different cell types: sieve cells (which conduct photosynthates), companion cells, phloem parenchyma, and phloem fibers. Unlike xylem conducting cells, phloem conducting cells are alive at maturity. The xylem and phloem always lie adjacent to each other. In stems, the xylem and the phloem form a vascular bundle. In roots, this is termed the vascular stele or vascular cylinder.

Adapted from Openstax biology 2e, section 25.1

 Core: Cell Biology

Western Blotting

JoVE 13381

Western blotting is an analytical technique for protein identification. It has various applications in immunology and medicine, including detecting diseases like bovine spongiform encephalopathy, mad cow disease, and human and feline immunodeficiency virus from biological samples.

The technique begins with separating proteins from the sample using sodium dodecyl sulfate-polyacrylamide gel electrophoresis (SDS-PAGE), followed by protein transfer, immunoblotting, and finally, protein detection.

The proteins from the polyacrylamide gel are transferred to specialized membranes like nitrocellulose (NC) or polyvinylidene difluoride (PVDF). Compared to NC, PVDF has higher durability and protein binding capacity and can be re-probed.

The protein transfer is either by electric current application (electroblotting) or capillary action. In electroblotting, the hydrophobic interactions transfer the negatively charged proteins from the gel and trap them on the membrane. In the other method, a filter paper placed on the gel soaks the buffer by capillary action, simultaneously drawing the proteins from the gel and transferring them to the membrane. To ensure the quality of transferred bands, non-specific protein stains like Ponceau S allow for a reversible check.

The next step, immunoblotting, involves treating the membrane with a blocking buffer like 3-5% Bovine serum albumin or non-fat dried milk. These buffers reduce the chance of non-specific binding of the primary antibody to the membrane. The primary antibodies can either be polyclonal or monoclonal and bind to the protein bands on the membrane through their Fab region. Next, the added secondary antibodies bind and attach to the Fc region of the primary antibodies. The secondary antibodies are tagged with enzymes which, upon suitable substrate addition, produce detectable colorimetric or chemiluminescent signals. Further, densitometry techniques can quantify the optical density of these imaged bands, corresponding to the abundance of proteins in the sample.

The expected result for a western blot can be affected by various factors like erroneous sample loading, inaccurate transfer with trapped air bubbles, and the use of non-specific antibodies. Also, inappropriate primary and secondary antibody concentrations and contaminated buffers can affect the quality of bands and their visualization.

 Core: Cell Biology

Atomic Force Microscopy

JoVE 13397

Atomic force microscopy (AFM) is a type of scanning probe microscopy that can analyze topographic details of various specimens like ceramics, glass, polymers, and biological samples. AFM offers over 1000 times more resolution than the optical imaging system. Images generated from AFM are three-dimensional surface profiles, offering an advantage over the flat, two-dimensional images from other imaging techniques.

The AFM Probe

The probe is regarded as the heart of any AFM setup and comprises the cantilever and tip assembly. Probes are the most commonly replaced part on this type of microscope because  the constant interaction with the samples wears down the tip. Therefore, the choice of material for the probe depends on the properties of the sample. Silicon probes, used to analyze hard samples, are stiffer and sharper than silicon nitride probes, which are better suited to scan softer samples. These sharp tips are produced using electrochemical etching or carbon nanotubes for higher accuracy analysis.

Imaging Modes of AFM

In AFM, surface topography is studied using the interaction between the probe tip and the sample surface. There are two main imaging modes — a static mode, also referred to as the contact mode, and a dynamic mode.

In the static or contact mode, the tip of the probe is in continuous contact with the sample surface. As the tip drags over the surface, repulsive forces between the sample and the tip result in the cantilever bending, which is recorded. The entire specimen surface is scanned back and forth in both x- and y-axes, called raster scanning, while the vertical movement of the cantilever records the z-axis, thus generating a 3D image.

In the dynamic mode, the probe oscillates just above the sample surface, coming close to, but not touching the surface. Attractive and repulsive forces determine the variation in distance between the tip and the sample, affecting the amplitude of cantilever oscillation. This feedback is recorded to construct the surface topography of the sample.

 Core: Cell Biology

Zygotic Development And Stem Cell Formation

JoVE 13461

The development of all multicellular organisms starts with the fusion of haploid cells called sperm and egg to form a diploid zygote. A zygote is a totipotent cell that can develop into a complete organism. The zygote undergoes cell division or cleavage to form an 8-cell mass. Until this stage, the cells are spherical, loosely attached, and remain totipotent. Totipotent cells are capable of developing both the embryonic and the extraembryonic tissues. However, as they continue to divide, they reach the 64-cell stage, called the blastocyst. The cells differentiate into two distinct developmental pathways and form the inner cell mass (ICM) and the trophectoderm (TE) cells. ICM is the source of embryonic stem cells or ES cells. ICM and TE are self-renewing but can differentiate only to limited cell types.

The ICM next differentiates into the embryonic tissues and forms the three germ layers (mesoderm, endoderm, and ectoderm) that give rise to all cell types of the embryo (neural cells, epithelial cells, muscle cells, and blood cells). In contrast, TE cells form the extraembryonic tissues such as the placenta, amnion, and chorion that cover the growing embryo. Once the blastocyst gets implanted on the uterine wall, they undergo further embryonic development.

After birth, each adult tissue retains a fraction of stem cells called the adult stem cells. Adult stem cells can form cells of a particular tissue type to help replace damaged or dead cells and maintain tissue integrity.  Adult stem cells include hematopoietic stem cells of the bone marrow and epidermal stem cells of the skin, gut, brain, lung, or liver.

 Core: Cell Biology

Embryonic Stem Cells

JoVE 13478

Embryonic stem (ES) cells were first discovered in mice in 1981 by Martin Evans. In 1998, James Thomson identified a method to isolate embryonic stem cells from humans. Human embryonic stem cells (hESCs) are obtained from 3-5 day old embryos that remain unused after an in vitro fertilization procedure.

ES cells are grown in a culture medium where they can divide indefinitely, creating ES cell lines. Under certain conditions, ES cells can differentiate, either spontaneously into a variety of cell types or in a directed fashion to produce desired cell types. Scientists can control which cell types are generated by manipulating the culture conditions—for example, changing the surface of the culture dish or adding specific growth factors to the culture medium—as well as by genetically modifying the cells. Researchers have generated many distinct cell types from ES cells, including blood, nerve, heart, bone, liver, and pancreatic cells through these methods.

Regenerative Medicine

Regenerative medicine concerns the creation of living, functional tissues to replace dead, diseased, or malfunctioning ones. ES cells are used in regenerative medicine because they can differentiate into any cell type. While this field is still in the early stages, several potentially beneficial cell types have been produced from ES cells, and clinical studies have begun to test their safety and effectiveness in patients. Some initial results have been promising. For instance, paralyzed patients regained some movement after receiving ES-derived nervous system cells. Additionally, ES cells can be used to study early events in human development and provide a source of specific cell types that can be used in drug testing and other scientific research.

Risks and Controversies

In hESC therapy, ES cells may cause immune rejection in recipient patients, as the embryonic cells are ‘foreign’ to the patient. Additionally, ES cells also carry the risk of tumor formation when implanted. Further, the use of hESCs is also controversial as it involves using human embryos. The controversy stems from the different views of an embryo held by scientists, religious leaders, and policy-makers.

 Core: Cell Biology

Faraday's Law

JoVE 13785

Faraday's law state that the induced emf is the negative change in the magnetic flux per unit of time. Any change in the magnetic field or change in the orientation of the area of the coil with respect to the magnetic field induces a voltage (emf). The magnetic flux measures the number of magnetic field lines through a given surface area. Magnetic flux is estimated from the integral of the dot product of the magnetic field vector and the area vector. The negative sign describes the direction in which the induced emf drives the current around a circuit. However, that direction is most easily determined with a rule known as Lenz's law.

In many practical applications, the circuit of interest consists of a number (N) of tightly wound turns. Each turn experiences the same magnetic flux. Therefore, the net magnetic flux through the circuit is N times the flux through one turn, and Faraday's law is written as the negative of N multiplied by the rate of change of magnetic flux through them.

Users of electrical appliances use a ground fault interrupter (GFI) device, which works as a circuit breaker, to avoid getting an electrical shock. This device works on the principle of Faraday's law. Another application of Faraday's law is seen in the electric guitar. In the electric guitar, a pick coil is placed near a vibrating string, and the coil is made of metal that can be magnetized. A permanent magnet inside the coil magnetizes the portion of the string while it is vibrating at a particular frequency, which results in a change in the magnetic flux. This change in flux produces an emf that is fed to the amplifier to produce a sound that listeners can hear.

 Core: Physics

Propagation of Uncertainty from Random Error

JoVE 14511

An experiment often consists of more than a single step. In this case, measurements at each step give rise to uncertainty. Because the measurements occur in successive steps, the uncertainty in one step necessarily contributes to that in the subsequent step. As we perform statistical analysis on these types of experiments, we must learn to account for the propagation of uncertainty from one step to the next. The propagation of uncertainty depends on the type of arithmetic operation performed on the values. For addition and subtraction, propagating the uncertainty requires us to express the absolute uncertainty of the outcome, which is the square root of the sum of absolute uncertainties for all steps. For multiplication and division, propagating uncertainty requires us to find the square root of the sum of the relative uncertainties for all steps, and this square root is equal to the relative uncertainty of the outcome, also known as the ratio between the absolute uncertainty of the outcome and the magnitude of the expected outcome. For exponential functions, we propagate the uncertainty by multiplying the power with the relative uncertainty of the base value, which then equates to the relative uncertainty of the outcome for the whole data set. Knowing how to propagate uncertainty correctly helps us identify the method that yields the least uncertainty, therefore optimizing our experimental protocols.

 Core: Analytical Chemistry

Chemical Equilibria: Redefining Equilibrium Constant

JoVE 14527

The effect of an inert salt on the solubility of a sparingly soluble salt is known as the salt effect. The degree of the salt effect varies with the ionic strength of the solution, which in turn depends on the activity of the species in the solution. The activity is expressed as the product of concentration and the activity coefficient of the species.

To calculate the equilibrium constants of solutions of moderately high ionic strength, one must account for the salt effect. This redefined equilibrium constant is also called the thermodynamic equilibrium constant or standard equilibrium constant, as it expresses the Gibbs energy change of the process. The thermodynamic equilibrium constant incorporates the ionic strength of the solution.

In solutions of low ionic strength (nearly an ideal solution), the activity coefficient is close to 1. Thus, the thermodynamic equilibrium constant is approximately equal to the concentration equilibrium constant.

The activity coefficient corrections are often ignored to simplify the experimental calculations of equilibrium constants. This approximation is valid for dilute solutions containing singly charged ions or non-dissociating species with ionic strengths lower than 0.01 mol/L. Activity coefficient corrections become more critical for solutions with ionic strengths greater than 0.01 mol/L or of multiply charged ions. Ignoring the activity coefficient in such cases results in significant errors in calculations.

 Core: Analytical Chemistry

Solution Composition During Acid/Base Titrations

JoVE 14543

The titration of a weak acid with a strong base results in the formation of water and the conjugate base of the acid. For instance, titrating acetic acid with sodium hydroxide leads to the formation of water and sodium acetate. A solution of acetic acid and sodium acetate constitutes a buffer whose relative concentration at different stages of the titration is indicated by the α values, which represent percentages of the weak acid and its conjugate base.

The α0 and α1 values represent the relative equilibrium concentration of acetic acid and acetate ions, respectively. Before adding the base, the α0 and α1 values are found to be 0.987 and 0.013, respectively. This suggests that the solution comprises 98.7% acetic acid and 1.3% acetate ions. With the addition of sodium hydroxide during the titration, α0 decreases while α1 increases. At the half-equivalence point, where the pH is equal to the pKa of acetic acid, the value of α0 is equal to α1. The solution at this stage is a mixture of 50% acetic acid and 50% acetate ions. As the titration proceeds, α0 decreases, and its value drops to zero at the equivalence point. On the contrary, the value of α1 becomes unity. 

 Core: Analytical Chemistry

Gravimetry: Overview

JoVE 14582

Gravimetric analysis is a quantitative method where the analyte is isolated and weighed directly or after conversion into a substance of known composition. Gravimetric analysis can be classified as precipitation, electrogravimetry, volatilization, and particulate gravimetry, based on the method used to isolate the analyte.

In precipitation gravimetry, the analyte is converted into a precipitate and weighed. For example, the silver content in a sample can be estimated by precipitating and weighing silver chloride. In electrogravimetry, the analyte undergoes a redox reaction in an electrolytic cell,  following which a solid film of the analyte is deposited over an electrode. Cu2+ and Pb2+ are ions that can be estimated by electrogravimetry. Volatilization gravimetry can be applied to a volatile analyte or an analyte that undergoes thermal/chemical decomposition to release a volatile component. The volatile species can be trapped and the mass determined by weighing the trapped species in an adsorbent trap. Alternatively, the loss of mass from the sample can be determined.

In particulate gravimetric analysis, particulate analytes are removed by filtration or extraction, and their mass is determined. This type of method is frequently employed in air and water quality monitoring systems.

 Core: Analytical Chemistry

Filtration

JoVE 14615

Filtration is a physical separation process that involves passing a suspension through a porous medium to separate solids from fluids. During filtration, solids collect on the porous medium while liquids, also collectively known as the filtrate, pass through. The filtration medium is selected based on the filtration purpose, quantity, and nature of the precipitate. The general criteria for a suitable filtering medium are that it is inert, mechanically strong, nonabsorbent toward dissolved materials and permissive toward rapid filtration.

The simplest filtration apparatus consists of filter paper fitted in a long-stemmed funnel sitting above a beaker. The solution to be filtered is poured down a glass rod onto the filter paper. The filtrate is collected in the beaker, and the solid is retained on the filter paper. Any solid that adheres to the glass rod or the beaker with the original solution is dislodged using a rubber policeman. Because filter paper is hygroscopic, ashless or low-ash filter paper is preferred when weighing is required. Ashless filter paper is also employed in gravimetric procedures that involve igniting the solids before weighing.

Alternatively, the solid can be collected in glass or silica crucibles containing a porous glass disc. The solid is transferred to the crucible fitted into a Buchner flask, and filtration is performed under suction. After filtration, the crucible is dried and weighed directly. The weight difference of the crucible before and after filtration gives the mass of the collected solid.

 Core: Analytical Chemistry

Excitation-contraction Coupling in Skeletal Muscles

JoVE 14842

Excitation-contraction coupling is a series of events that occur between generating an action potential and initiating a muscle contraction. It occurs at the triad, a structure found in skeletal muscle fibers that comprise a T-tubule and terminal cisternae of the sarcoplasmic reticulum on each side. These triads are visible in longitudinally sectioned muscle fibers. They are typically located at the A-I junction — the junction between the A and I bands of the sarcomere.

When an action potential reaches a triad, it triggers calcium release from the sarcoplasmic reticulum. This temporary change in calcium permeability lasts for about 0.03 seconds, significantly increasing calcium concentration in and around the sarcomere. Troponin acts as a lock, preventing interaction between thick and thin filaments within a sarcomere, the smallest functional unit of muscle cells. Upon calcium binding, troponin undergoes a conformational change, allowing tropomyosin to move away from the active sites on actin filaments. This change marks the beginning of the contraction cycle.

Once the active sites are exposed, the myosin heads attach to them, creating cross-bridges between thick and thin filaments. The connection between the head and tail of the myosin molecule acts as a hinge, allowing the heads to pivot. This pivoting motion, fueled by the energy released from ATP hydrolysis, is called the power stroke and is a crucial step in muscle contraction. The myosin heads slide the thin filaments from both ends of a sarcomere toward the M line. This inward movement causes the filaments to move towards the center of the sarcomere. The filaments may often overlap, causing the I band and H zone to narrow and disappear when the muscle fully contracts. However, the A band and lengths of the thick and thin filaments do not change. Since the thin filaments at each side of the sarcomere attach to Z discs, their inward movement results in the shortening of the sarcomere and, as a result, the entire muscle fiber.

 Core: Anatomy and Physiology

Naming Skeletal Muscles

JoVE 14864

The naming of the approximately 700 muscles in the human body is based on a set of criteria designed to provide descriptive information about each muscle, making it easier to identify and remember them.

The key factors used in naming muscles include:

  1. Region of the Body: Certain muscles in the body are classified based on their location or association with a specific area. For example, the temporalis muscle is in the head region, whereas the brachialis is in the arm.
  2. Position and Depth: Muscles visible on the body's surface are often termed 'externus' or 'superficialis', while deeper muscles are referred to as 'internus' or 'profundus'. Muscles that position or stabilize an organ are called 'extrinsic', while those located entirely within an organ are 'intrinsic'.
  3. Direction: Some muscles are named according to the direction in which they run relative to the body's longitudinal axis. 'Transversus' refers to muscles running across this axis, while 'oblique' denotes muscles running at a slanting angle.
  4. Fascicle Arrangement: The term 'rectus' is used for muscles with straight fascicles running along their longitudinal axis. For example, 'rectus abdominis' refers to a muscle in the abdomen with straight fascicles.
  5. Structural Characteristics: Some muscles are named after their distinctive structural features, such as the number of tendons of origin, shape, and size. For example, 'biceps brachii' has two tendons of origin, 'trapezius' is trapezoid-shaped, and 'longus' refers to long muscles.
  6. Origin and Insertion: Many muscle names indicate the specific origin and insertion points. For example, the 'genioglossus' originates at the chin or 'geneion' and inserts in the tongue or 'glossus'.
  7. Action: Many muscles are named after their primary function or action, such as flexor, extensor, pronator, abductor, adductor, and rotator.
  8. Particular Occupations or Habits: A few muscles carry names related to specific movements associated with certain occupations or habits. For instance, the 'buccinator' compresses the cheeks as if blowing a trumpet, and the 'sartorius', the longest muscle in the body, is active when crossing legs, reminiscent of a tailor's sitting position.

 Core: Anatomy and Physiology

Muscles of the Leg that Move the Foot and Toes

JoVE 14881

The human leg comprises an intricate system of muscles that facilitate the movement of feet and toes. Within this system, the muscles are categorized into the anterior, lateral, and posterior compartments, each with a unique set of muscles carrying out specific functions.

Anterior Compartment

The anterior compartment includes muscles that contribute to the dorsiflexion of the foot. This compartment houses the tibialis anterior, extensor hallucis longus, and extensor digitorum longus muscles. The tibialis anterior muscle, situated along the tibia's lateral surface, is easily felt and is the thickest. The tibialis anterior also causes foot inversion. Additionally, the extensor hallucis longus and extensor digitorum longus extend the great toe and other toes, respectively. The fibularis tertius muscle, originating from the same source as the extensor digitorum longus, is also present here. It everts the foot.

Shin splints, medically referred to as medial tibial stress syndrome, are a common affliction among athletes, especially runners and dancers. It results from overuse and repetitive stress of the anterior compartment muscles. Preventive measures include regular stretching, strength training, and choosing appropriate footwear.

Lateral Compartment

The lateral compartment contains the fibularis longus and fibularis brevis muscles. The fibularis longus is the longer and more superficial of the two muscles. It originates from the head and upper portion of the fibula, the outer and thinner of the two bones in the lower leg, and the adjacent intermuscular septa. The fibularis brevis, lying just beneath the fibularis longus, is shorter and originates from the lower two-thirds of the fibula. These two muscles help in the plantar flexion and eversion of the foot.

Posterior Compartment

The posterior compartment is further divided into superficial and deep muscle groups. The superficial group includes the gastrocnemius, soleus, and plantaris muscles, commonly called the calf muscles. These muscles share the Achilles tendon, the body's most robust tendon, which attaches to the calcaneal bone of the ankle. The human upright stance is attributed to the size of these muscles. The gastrocnemius forms the calf's prominent part, while the soleus, found beneath the gastrocnemius, bears a flat, broad look. The plantaris, a small muscle, is nestled between these two muscles and is absent in some individuals.

The deep posterior compartment houses the popliteus, tibialis posterior, flexor digitorum longus, and flexor hallucis longus muscles. The popliteus forms the floor of the popliteal fossa. The tibialis posterior is the deepest muscle in the posterior compartment, positioned between the flexor digitorum longus and flexor hallucis longus muscles. Despite flexing just the big toe at the interphalangeal joint, the flexor hallucis longus is larger than the flexor digitorum longus, which flexes the other four toes.

 Core: Anatomy and Physiology

Postsynaptic Potential (PSP)

JoVE 14898

Postsynaptic potential (PSP) refers to a change in the electrical potential of a neuron when neurotransmitters released by presynaptic neurons bind to postsynaptic receptors. This potential can either be excitatory, leading to depolarization and ultimately action potential generation, or inhibitory, leading to hyperpolarization and suppression of the postsynaptic neuron.

There are two types of receptors: ionotropic and metabotropic.

The ionotropic receptor is the membrane protein that has an ion channel and a neurotransmitter binding site to facilitate the ion flow in or out of the postsynaptic neuron and alter its membrane potential. A metabotropic receptor has a binding site but lacks the ion channel and acts through G protein.

The PSP generated by ionotropic receptors is either excitatory or inhibitory, depending on the type of ion moving through the channel. For example, an influx of positive ions (typically sodium) will cause depolarization and create an excitatory PSP, while an efflux of positive ions (like potassium) will produce a hyperpolarizing response due to increased negativity in the neuron, forming an inhibitory PSP.

An example of excitatory postsynaptic potential (EPSP) is found in the synapses of neurons in the central nervous system, where neurotransmitters such as acetylcholine, glutamate, and GABA are released to cause depolarization of the postsynaptic membrane and, ultimately, an action potential that propagates down the axon to its target cell. Inhibitory postsynaptic potentials (IPSPs) occur when inhibitory neurotransmitters such as glycine or GABA are released from a presynaptic neuron and cause hyperpolarization on the postsynaptic neuron, making it less likely for an action potential to be generated.

The PSP plays an important role in the functioning of neural circuits, allowing neurons to communicate with each other and ultimately coordinating the activity of a large network of neurons. PSPs help shape and modulate the output of a neuron in response to its inputs and are important for learning, memory formation, and information processing. In addition, PSPs are used as a readout measure in neuroscience studies, providing valuable insights into neural function. PSPs are essential for intercellular communication within the nervous system, allowing it to process and respond to external stimuli in an appropriate manner. By studying how postsynaptic potentials work at the cellular level, we can gain a better understanding of how neuronal networks operate to form complex behaviors.

Postsynaptic potentials are essential for information processing in the nervous system, and thus, it is important to understand how they work in order to better understand normal neurological functioning and diseases such as epilepsy or Parkinson's disease. By studying postsynaptic potentials, researchers can get a better insight into how neurons communicate with each other, allowing us to develop more effective treatments for these conditions.

 Core: Anatomy and Physiology

Brainstem: Control Centers of Medulla

JoVE 14914

The medulla oblongata is a crucial part of the brainstem responsible for controlling various autonomic and involuntary functions. It contains several nuclei, including the olivary, cuneate, gracile, and solitary nuclei.

Olivary Nucleus

The olivary nucleus, or inferior olivary nucleus, is located within the ventrolateral part of the medulla oblongata. It is primarily involved in motor coordination and motor learning. The olivary nucleus receives input from the spinal cord, cerebellum, and motor cortex and sends output to the cerebellum via climbing fibers. This communication helps fine-tune motor movements and maintain balance.

Cuneate Nucleus

The cuneate nucleus is found laterally in the medulla oblongata and receives input from the cuneate fasciculus, which carries sensory information from the upper body's proprioceptors. It is involved in processing touch, pressure, vibration, and proprioceptive sensations. The cuneate nucleus then sends this information to the thalamus via the medial lemniscus pathway, contributing to motor control of skeletal muscles.

Gracile Nucleus

The gracile nucleus is situated medially in the medulla oblongata. It receives input from the gracile fasciculus, which carries sensory information from the lower body's proprioceptors. Like the cuneate nucleus, the gracile nucleus processes touch, pressure, vibration, and proprioceptive sensations and sends this information to the thalamus via the medial lemniscus pathway.

Solitary Nucleus

The solitary nucleus, also known as the nucleus tractus solitarius (NTS), is located in the dorsal part of the medulla oblongata. It receives sensory input from the facial, glossopharyngeal, and vagus nerves, which are responsible for taste sensation and visceral information from the thoracic and abdominal organs. The solitary nucleus regulates cardiovascular and respiratory activities, gastrointestinal motility, and taste sensation.

Control Centers

The medulla oblongata houses several control centers vital for maintaining bodily functions. The cardiovascular center regulates heart rate and blood vessel diameter, thereby controlling blood pressure. The respiratory center controls the rate and depth of breathing by receiving input from chemoreceptors and mechanoreceptors. It also harbors nuclei responsible for initiating reflex actions such as sneezing, hiccupping, and coughing to protect the respiratory system.

The solitary nucleus serves as the control center that regulates gastrointestinal motility by processing visceral information from the thoracic and abdominal organs. The solitary nucleus also processes taste sensations by receiving sensory input from the facial, glossopharyngeal, and vagus nerves. Both the cuneate and gracile nuclei contribute to motor control by processing sensory information from the upper and lower body's proprioceptors, respectively.

 Core: Anatomy and Physiology

Introduction to Sensory Receptors

JoVE 14932

Sensory receptors are vital in our ability to perceive and interpret the world. Sensory receptors are specialized cells in the peripheral nervous system that respond to various stimuli and enable one to experience different sensations. Based on specific criteria, sensory receptors are classified into distinct types.

The first classification criterion is based on cell type, position, and function. Some receptor cells are neurons with free nerve endings, where their dendrites are embedded in the tissue that receives the sensation. Others are neurons with encapsulated endings, surrounded by connective tissue, enhancing their sensitivity to stimuli. The third category of cells consists of specialized receptor cells, which are equipped with unique structural elements designed to detect and interpret specific kinds of stimuli. For example, photoreceptors are located in the eyes, while gustatory receptors can be found within the taste buds.

Receptors can also be classified based on their location relative to the stimuli. Exteroceptors, such as somatosensory receptors in the skin, are located near stimuli in the external environment, enabling us to perceive touch and pressure. Conversely, interoceptors interpret stimuli from internal organs and tissues, like receptors that sense changes in blood pressure in the aorta or carotid sinus. Finally, proprioceptors are located near moving body parts, such as muscles, and help us interpret body positions as they move, contributing to our sense of balance and coordination.

Physical stimuli, such as pressure, vibration, sound, and body position, are interpreted through mechanoreceptors. Temperature, a critical sensory aspect, is sensed through thermoreceptors sensitive to temperatures above or below normal body temperature.

 Core: Anatomy and Physiology

Sympathetic Pathways: Sympathetic Chain Ganglia

JoVE 14949

The sympathetic chain ganglia, also known as the sympathetic trunk ganglia or paravertebral ganglia, are a series of ganglia located bilaterally on either side of the spinal column. These ganglia serve as relay stations for the sympathetic nervous system. Preganglionic neurons originating in the spinal cord project their axons to the sympathetic chain ganglia. Within the ganglia, these preganglionic fibers synapse with postganglionic neurons.

The postganglionic neurons of the sympathetic trunk ganglia follow four pathways that leave the ganglia to innervate their target effector organs.

  1. Spinal nerves: Some postganglionic axons merge with spinal nerves through the gray rami communicans, providing sympathetic innervation to the skin of the neck, limbs, and torso.
  2. Cephalic periarterial nerves: Some postganglionic axons wrap around arteries, like the carotid arteries, supplying sympathetic innervation to the face, including sweat glands, blood vessels, hair follicles, and various structures in the head.
  3. Sympathetic nerves: Some postganglionic axons leave the sympathetic trunk, forming sympathetic nerves that reach visceral effectors in the thoracic cavity, such as the heart and lungs.
  4. Splanchnic nerves: Some preganglionic fibers pass through the sympathetic trunk without synapsing, forming splanchnic nerves that travel to the collateral ganglia.

 Core: Anatomy and Physiology

The Physiology of Taste

JoVE 14968

The perception of a salty flavor is facilitated by sodium ions within the oral salivary fluid. Upon consumption of a salty substance, salt crystals disassemble, leading to the liberation of its constituents—Na+ and Cl- ions. These ions subsequently dissolve into the salivary fluid present in the oral cavity. The external environment of the gustatory cells experiences an elevation in Na+ concentration, thereby establishing a potent concentration gradient. This gradient propels the diffusion of Na+ ions into these cells. The influx of Na+ triggers the phenomenon of cell membrane depolarization, subsequently evoking a receptor potential.

Perception of sourness is associated with the detection of hydrogen ion concentration. Analogous to the role of sodium ions in evoking saltiness, hydrogen ions permeate the cellular membrane, resulting in depolarization. Sourness is a tactile response to acids present in our edibles. An increased hydrogen ion concentration in salivary fluid, corresponding to decreased salivary pH, elicits graded potentials within gustatory cells. For instance, citric acid-laden orange juice manifests a sour taste due to its pH value approximating 3. However, it's often sweetened to obscure the inherent sourness.

The salty and sour tastes are induced by cations such as Na+ and H+. The remaining tastes result from food molecules contacting a specific receptor type, a G protein-coupled receptor. This interaction activates a G protein signaling pathway, culminating in the depolarization of the gustatory cell. Sweetness is perceived when gustatory cells detect glucose molecules dissolved in saliva. However, other monosaccharides, such as fructose and artificial sweeteners, including aspartame, saccharin, or sucralose, also stimulate sweet receptors. Each of these compounds has a different binding affinity to the G protein–coupled receptor, which is why some may be perceived as sweeter than glucose.

The bitter taste sensation, akin to sweetness, is triggered when food molecules attach to G protein-coupled receptors. However, the underlying mechanisms vary significantly due to the broad spectrum of bitter-flavored compounds. Some of these substances depolarize or hyperpolarize gustatory cells, whereas others modulate G protein activation within these cells. The specific response elicited is contingent on the molecular constitution of the receptor-bound compound. A prominent class of bitter compounds is represented by alkaloids, nitrogen-rich substances ubiquitously found in plant products like coffee, hops, tannins, tea, and medications such as aspirin. These toxic alkaloids render the plant less prone to microbial invasion and less appealing to herbivorous organisms, suggesting that the function of bitter taste may be principally linked to the activation of protective reflexes, such as the gag reflex, to prevent the ingestion of potential toxins. This means that traditionally consumed bitter foods are usually paired with sweet components to render them palatable (for instance, adding cream and sugar to coffee). Notably, the posterior region of the tongue, possessing the highest concentration of bitter receptors, is an effective site for triggering the gag reflex, providing a mechanism to expel potentially toxic substances.

Umami, frequently described by its savory flavor, is akin to the sweet and bitter tastes and originates from stimulating G-protein-linked receptors by a distinct molecule. This essential molecule, L-glutamate, an amino acid, is the initiator of this receptor. As a result, the umami sensation is frequently experienced when consuming foods rich in protein. Consequently, it's not unexpected that meals containing a high proportion of meat carry a savory descriptor.

Upon activation by taste molecules, gustatory cells initiate a release of neurotransmitters. These neurotransmitters subsequently interact with the dendrites of sensory neurons. Included within these neurons are components of the facial and glossopharyngeal cranial nerves, as well as a segment of the vagus nerve dedicated to the gag reflex. Specifically, the facial nerve connects with taste buds in the tongue's anterior third. In contrast, the glossopharyngeal nerve links with taste buds in the posterior two-thirds of the tongue. Lastly, the vagus nerve communicates with taste buds near the far posterior of the tongue, bordering the pharynx, which showcases heightened sensitivity to harmful stimuli, such as bitterness.

 Core: Anatomy and Physiology

Synthesis and Functions of Calcitonin

JoVE 14984

Calcitonin, a vital polypeptide hormone, regulates calcium levels within body fluids. It is released by the parafollicular cells, also known as C cells, situated in the follicular epithelium of the thyroid gland. Calcitonin responds to fluctuations in blood calcium levels and the influence of gastrointestinal hormones like gastrin and cholecystokinin.

The exact mechanisms by which calcitonin operates in calcium homeostasis remain elusive, but its significance is evident in several vital functions. Notably, calcitonin is crucial in bone calcium storage, particularly during childhood, stimulating bone growth and mineral deposition. Additionally, calcitonin prevents bone mass loss in adults, especially in prolonged starvation and later stages of pregnancy.

Calcitonin's impact extends to the digestive system, which participates in calcium ion absorption. Additionally, it exerts a regulatory influence on post-meal calcium surges by promoting calcium excretion through the kidneys.

The intricate involvement of calcitonin in calcium regulation underscores its dual role in maintaining skeletal integrity and contributing to overall calcium homeostasis. This hormone regulates calcium balance in the body throughout different life stages and physiological conditions.

 Core: Anatomy and Physiology

Complexometric Titration: Ligands

JoVE 15002

Different monodentate and polydentate ligands are used as complexing agents in complexometric titration reactions. The formation of complexes by mono- and bidentate ligands involves two or more intermediate steps, limiting their use as complexing agents. In comparison, polydentate ligands can form complexes with metal ions in a single-step process, facilitating sharper end points. This means polydentate ligands, such as amino carboxylic acid derivatives, are most commonly employed in complexometric titrations. These ligands can form stable chelates with metal ions via carboxylate oxygens and amine nitrogens. Among the amino carboxylic acid derivatives, ethylenediaminetetraacetic acid (EDTA) is the most widely used complexing agent.

Complexometric titrations are used to quantify most metals. Quantitative determination of the total hardness of water and the investigation of calcium in the blood are two such applications of complexometric titration.

 Core: Analytical Chemistry

Ohm's Law

JoVE 15069

Resistors are fundamental components in electrical circuits, often manufactured from metallic alloys or carbon compounds. They model a material's ability to resist the flow of electric current, a characteristic that is crucial in controlling and regulating electrical power within a circuit.

This current-resisting behavior of resistors is governed by Ohm's law, which states that the voltage across a resistor is directly proportional to the current flowing through it.

Equation1

Equation2

The constant of proportionality in this relationship is known as resistance. Measured in units called ohms (Ω), resistance indicates how effectively a material can impede the flow of electric current.

The resistance of a material is directly proportional to its length and inversely proportional to its cross-sectional area. The proportionality constant in this relationship is known as resistivity, which varies among different materials. Conductors, which allow easy passage of electrical current, have low resistivities, while insulators, which block the flow of current, have high resistivities.

Resistors can further be categorized into linear and non-linear based on their voltage-current relationship. Linear resistors obey Ohm's law, exhibiting a linear relationship between their current and voltage. Non-linear resistors, on the other hand, do not adhere to this law, and their resistance changes with the applied voltage or current.

Two extreme cases of resistance are short circuits and open circuits. A short circuit has zero resistance and voltage drop, allowing it to carry any amount of current. In contrast, an open circuit has infinite resistance, resulting in zero current flow, and can sustain an unrestricted voltage.

Conversely, the reciprocal of resistance is known as conductance. This measures how well an element conducts electric current. Conductance is measured in Siemens (S), with one Siemens being equivalent to one ampere per volt.

 Core: Electrical Engineering

Energy Stored in Capacitors

JoVE 15085

A parallel plate capacitor, when connected to a battery, develops a potential difference across its plates. This potential difference is key to the operation of the capacitor, as it determines how much electrical energy the capacitor can store.

By integrating the equation that relates voltage and current in a capacitor, one can derive an equation for the voltage across the capacitor at any given time. This equation is crucial in understanding and predicting the behavior of capacitors in electronic circuits.

One interesting property of capacitors is that they possess a kind of memory. The voltage across a capacitor at any moment depends on the past flow of current through it. This means that capacitors can "remember" their charging and discharging history, which can be useful in various applications such as memory storage in computers.

The instantaneous power delivered to a capacitor can be used to determine the amount of energy stored in the capacitor. If we consider an uncharged capacitor at time equals minus infinity, it has zero voltage. This means that the energy stored in the capacitor can be determined in terms of charge and capacitance. This represents the energy present in the electric field between the plates.

This stored energy can be retrieved in terms of power since an ideal capacitor does not dissipate energy. However, real-world capacitors are not ideal. A non-ideal capacitor has a parallel-model leakage resistance, but this is usually high enough to be neglected in most practical applications.

A unique characteristic of capacitors is that they act as an open circuit to direct current (DC) voltage but can get charged when connected to a battery. This property allows capacitors to block DC while letting alternating current (AC) pass.

Another important feature of capacitors is that the voltage across them is always continuous and cannot change abruptly. This behavior is essential in many applications, such as smoothing out voltage in power supplies and filtering out noise in signal processing.

In conclusion, understanding capacitors, their properties, and their behavior in circuits is a fundamental aspect of electronics. It enables engineers to design and construct complex electronic systems that are integral to modern life.

 Core: Electrical Engineering

Second-order Op Amp Circuits

JoVE 15103

Implementing second-order low-pass filters in audio systems is crucial in refining audio signals by eliminating undesirable high-frequency noise. These filters typically involve second-order op-amp circuits configured as voltage followers, encompassing two nodes with distinct storage elements.

The analysis of such circuits follows a systematic approach, similar to the second-order RLC circuits. In practical scenarios, bulky inductors are rarely employed due to their size and weight. This means the focus here is primarily on RC second-order op-amp circuits, which have extensive applications in devices like filters and oscillators.

Two differential equations emerge after applying Kirchhoff's current law at the nodes. Furthermore, a second-order characteristic differential equation is deduced by observing voltage relationships across the circuit components. This equation embodies both transient and steady-state responses.

Equation1

The transient response gradually diminishes over time and shares resemblances with the solutions found in source-free circuits, exhibiting characteristics of overdamped, underdamped, and critically damped scenarios. Once the circuit achieves a steady state, the capacitors and resistors no longer conduct current, and the ideal op-amp input terminals block current flow. Consequently, the steady-state response matches the source voltage.

Notably, eliminating the input source voltage leads to a pure transient response. These second-order op-amp circuits have diverse applications in enhancing audio quality and are pivotal in various audio processing systems. Understanding their behavior under different damping scenarios aids in achieving optimal audio signal refinement.

 Core: Electrical Engineering

Design Example: Vintage Mixing Console

JoVE 15172

A sound engineer at a music company recently encountered a problem. The output from their newly acquired studio's vintage mixing console was too low for the requirements of modern recording equipment. To rectify this situation, the engineer decided to design an audio pre-amplifier using an operational amplifier (op-amp) to boost the signal level.

The specifications for the pre-amplifier were clear. It needed to amplify the audio signal by a factor of 10, have an input impedance above 10 kilo-ohms to prevent overloading the console, and be capable of covering the entire audio spectrum from 20 Hz to 20 kHz.

The engineer chose the op-amp model 741, a popular choice in audio applications due to its ability to handle the required frequency range without significant gain loss. This model would ensure that the pre-amplifier could adequately amplify the signals across the whole audio spectrum.

To meet the requirement of a gain of 10 and high input impedance, the engineer decided to use the non-inverting operational amplifier configuration. The gain of a non-inverting amplifier equals one plus the ratio of the feedback resistor to the resistor connected to the inverting input.

To achieve the desired gain of 10, the engineer selected a 100 kilo-ohms feedback resistor and a 10 kilo-ohms input resistor. This combination would provide the necessary amplification while maintaining an adequate input impedance to prevent overloading the mixing console.

This carefully designed pre-amplifier effectively boosted the console's output, enabling it to drive the digital recording equipment successfully. The result was a seamless integration of vintage and modern equipment, ensuring the preservation of the unique sound qualities of the vintage console while taking advantage of the capabilities of contemporary digital recording technology.

 Core: Electrical Engineering

An Introduction to the Zebrafish: Danio rerio

JoVE 5128

Zebrafish (Danio rerio) are small freshwater fish that are used as model organisms for biomedical research. The many strengths of these fish include their high degree of genetic conservation with humans and their simple, inexpensive maintenance. Additionally, gene expression can be easily manipulated in zebrafish embryos, and their transparency allows for observation of developmental processes.

This overview video first introduces basic zebrafish biology, including their phylogeny, life cycle, and natural environment, before presenting the features that make them so useful in the lab. A brief history of zebrafish research is also provided through a review of major discoveries made in fish, ranging from the early establishment of methods for efficient genetic screening to the discovery of novel therapeutics for human diseases such as cancer. Finally, some of the many avenues of experimentation performed in zebrafish are discussed, including immunological and developmental studies.

 Biology II

Invertebrate Lifespan Quantification

JoVE 5338

Many animals naturally stop growing upon reaching adulthood, after which they undergo aging or "senescence" until dying. The amount of time between an organism\'s birth and death is called its lifespan, which can be influenced by various biological and environmental factors. By exposing organisms to different growth conditions, scientists can better understand the factors affecting lifespan. Flies and worms are ideal organisms to perform such experiments, given their short generation time and simple culture requirements.

This video provides a brief overview of the factors affecting aging, and goes on to describe basic protocols for invertebrate lifespan quantification experiments. Finally, three research applications of lifespan quantification will be discussed. These experiments explore the effects of diverse factors, such as temperature, drugs, pathogens, and diet, on lifespan.

 Developmental Biology

Modeling Social Stress

JoVE 5429

Stress negatively affects our quality of life. In particular, some individuals experience social stress when placed in a social environment that they are unfamiliar with or have difficulty adjusting to. Since it is hard to examine mechanisms of social stress in humans, modeling this condition in animals may help scientist in developing new therapies for treating this commonly encountered problem.

This science education video begins by discussing the known anatomy and physiology behind stress response. Then, we explain a well-established paradigm for modeling social stress in rodents, the Resident-Intruder task. In the applications section, we review some example studies in which response to stress is measured.

 Behavioral Science

SNP Genotyping

JoVE 5544

Single nucleotide polymorphisms, or SNPs, are the most common form of genetic variation in humans. These differences at individual bases in the DNA often do not directly affect gene expression, but in many cases can still be useful for locating disease-associated genes or for diagnosing patients. Numerous methodologies have been established to identify, or “genotype”, SNPs.

JoVE’s introduction to SNP Genotyping begins by discussing what SNPs are and how they can be used to identify disease-associated genes. Several common SNP genotyping methods are then examined, including direct hybridization, PCR-based methods, fragment analysis, and sequencing. Finally, we present several examples of how these techniques are applied to genetic research today.

 Genetics

Electrochemical Measurements of Supported Catalysts Using a Potentiostat/Galvanostat

JoVE 5698

Source: Laboratory of Dr. Yuriy Román — Massachusetts Institute of Technology

A potentiostat/galvanostat (often referred to as simply a potentiostat) is an instrument that measures current at an applied potential (potentiostatic operation) or measures potential at an applied current (galvanostatic operation) (Figure 1). It is the most commonly used instrument in the electrochemical characterization of anode and cathode materials for fuel cells, electrolyzers, batteries, and supercapacitors.

Conventionally, these anode and cathode materials are interfaced with a potentiostat via a three-electrode electrochemical cell. The electrode leads from the potentiostat are connected to the reference electrode, the counter electrode (often called the auxiliary electrode), and the working electrode (which contains the test material of interest). The electrochemical cell is then filled with a high ionic strength electrolyte solution, such as an acidic, alkaline, or salt solution. The media for this high ionic strength solution is typically aqueous; however, for applications necessitating higher operating cell potential windows, such as batteries and supercapacitors, non-aqueous media is often used. The cell media is degassed with an inert gas (to prevent unwanted side reactions) or with a test gas (if the test reaction involves a gas at one of the electrodes).

Alternatively, a salt bridge or membrane is employed to maintain ionic contact if the two half cells are to be measured in different electrolytes. In heterogeneous electrocatalysis, this type of "two compartment" cell is often used if the test molecule at the working electrode is also reactive at the counterelectrode. This happens frequently as the counterelectrode typically employed is platinum, which is a highly active catalyst for many reactions. Here, single compartment cells will be used, where all three electrodes are in the same media.

This video will explain the process of polishing a working electrode, preparing a catalyst ink, mounting the catalyst ink onto the working electrode, preparing the electrochemical cell, and then performing electrochemical measurements. The measurements that are performed include: cyclic voltammetry (CV), linear sweep voltammetry (LSV), chronopotentiometry (CP), and chronoamperometry (CA).

Figure 1
Figure 1. An example of a single compartment electrochemical cell. a.) Teflon cap, b.) glass cell, c.) Pt wire counter electrode, d.) working electrode, e.) Ag/AgCl reference electrode, f.) 0.5 M aqueous sulfuric acid electrolyte solution.

 Analytical Chemistry

Overview of Biomaterials

JoVE 5797

Biomaterials are materials engineered to interact favorably with biological organisms or molecules. These materials can be derived from or produced by an organism, or can even be a synthesized polymer. Engineers use these novel materials in a wide range of applications, such as tissue engineering, biosensing and drug delivery.

This video introduces common biologically derived materials, and provides examples of common techniques used to process them. Key challenges in the field are discussed, along with several applications of these methods.

 Bioengineering

Chair Conformation of Cyclohexane

JoVE 11713

The chair conformation is the most stable form of cyclohexane due to the absence of angle and torsional strain. The absence of angle strain is a result of cyclohexane’s bond angle being very close to the ideal tetrahedral bond angle of 109.5° in its chair conformer. Similarly, the torsional strain is also absent owing to the perfectly staggered arrangement of bonds.

The hydrogen atoms linked to carbons are arranged in two different axial and equatorial orientations to achieve this staggered form. The axial bonds are directed straight up or down, lying parallel to the ring axis, whereas the equatorial bonds are pointed sideways roughly along the equator of the ring. Out of the six axial bonds, three are pointed up, and the remaining three are pointed downward. Similarly, three bonds are slanted upwards among the six equatorial bonds, while the remaining three are slanted downwards. Thus, each carbon atom in the cyclohexane ring has an axial and an equatorial bond, pointing in opposite directions.

A chair conformation of cyclohexane can undergo a conformational change into another chair conformer by the partial rotation of C-C bonds. This chair-chair interconversion that leads to the generation of two equivalent energy forms is known as ring flipping. Upon ring flipping, the axial and equatorial bonds interchange their positions. The axial bonds in one chair conformation get converted to equatorial bonds in the other chair conformation, while equatorial bonds change their position to axial bonds.

 Core: Organic Chemistry

Nomenclature of Alkenes

JoVE 11765

The IUPAC naming system for alkenes replaces -an- with -en- in the corresponding parent alkanes. Accordingly, a simple alkene replaces the -ane suffix of the alkane with -ene.

As per the IUPAC rules, the longest carbon chain containing the maximum number of double bonds is identified as the parent chain and is numbered such that the doubly bonded carbon atoms receive the lowest possible numbers. The location of the double bond is indicated by the number of its first carbon atom. In branched alkenes, the preference of numbering is given to the double bond over the substituents. The substituents are numbered according to their position in the parent chain and are listed alphabetically. However, the name and position of the substituent groups are cited before the double bond. In cycloalkenes, the ring is considered the parent chain, and the lowest possible numbers are assigned to the double bond. Alkenes having multiple double bonds in the parent chain are termed polyenes and have the infixes -adien- and -atrien- for two and three double bonds, respectively.

Some smaller alkenes have IUPAC-acknowledged common names. For example, ethene is commonly known as ethylene, propene is known as propylene, and so on.

 Core: Organic Chemistry

Oxidative Cleavage of Alkenes: Ozonolysis

JoVE 11784

In ozonolysis, ozone is used to cleave a carbon–carbon double bond to form aldehydes and ketones, or carboxylic acids, depending on the work-up.

Ozone is a symmetrical bent molecule stabilized by a resonance structure.

Figure1

Ozonolysis proceeds through an oxidative cleavage reaction. The first step is the electrophilic addition of ozone across the alkene double bond, forming an unstable molozonide intermediate, which reacts further to form a carbonyl and a carbonyl oxide. These intermediates rearrange to form an ozonide.

Figure2

The ozonide is treated with a mild reducing agent such as dimethyl sulfide or zinc to yield the carbonyl compounds as the final product.

Figure3

Ozonolysis with Different Substituted Alkenes

The conversion of ozonide to aldehydes, ketones, or carboxylic acids depends on the structure of the alkene starting material and different reaction conditions.

When a reductive work-up is used,  ozonolysis of monosubstituted alkenes such as 1-butene yields a mixture of aldehydes.

Figure4

Trisubstituted alkenes, such as 2-methyl-2-butene, on the other hand, form an aldehyde and a ketone.

Figure5

When an oxidative work-up is used, the reaction yields a ketone and an aldehyde that is further oxidized to the corresponding carboxylic acid.

Figure6

 Core: Organic Chemistry

Reduction of Alkynes to trans-Alkenes: Sodium in Liquid Ammonia

JoVE 11842

Alkynes can be reduced to trans-alkenes using sodium or lithium in liquid ammonia. The reaction, known as dissolving metal reduction, proceeds with an anti addition of hydrogen across the carbon–carbon triple bond to form the trans product. Since ammonia exists as a gas (bp = −33°C) at room temperature, the reaction is carried out at low temperatures using a mixture of dry ice (sublimes at −78°C) and acetone. 

When dissolved in liquid ammonia, an alkali metal, such as sodium, dissociates into a cation and a free electron. Ammonia molecules surround the free electrons, creating solvated electrons that impart a blue color to the solution. Solvated electrons are strong reducing agents and readily add to the alkyne triple bond.

Figure1

Figure2

Limitation:

The reduction of terminal alkynes with sodium in liquid ammonia does not proceed as efficiently as the reduction of internal alkynes. This is because terminal alkynes have acidic protons that readily react with the sodium–liquid ammonia mixture to form sodium acetylide. Stoichiometrically, three moles of a terminal alkyne undergo metal-dissolved reduction to give only one mole of the corresponding alkene and two moles of sodium acetylide.

Figure3

Therefore, the reaction conditions need to be modified to completely convert terminal alkynes to alkenes. A common approach involves adding ammonium sulfate to the reaction mixture. The ammonium ion released into the solution protonates the acetylide, thus preserving the terminal alkyne for subsequent reduction.

 Core: Organic Chemistry

Acid-Catalyzed Dehydration of Alcohols to Alkenes

JoVE 11924

In a dehydration reaction, a hydroxyl group in an alcohol is eliminated along with the hydrogen from an adjacent carbon. Here, the products are an alkene and a molecule of water. Dehydration of alcohols is generally achieved by heating in the presence of an acid catalyst. While the dehydration of primary alcohols requires high temperatures and acid concentrations, secondary and tertiary alcohols can lose a water molecule under relatively mild conditions.

Figure1

The acid-catalyzed dehydration of secondary and tertiary alcohols proceeds via an E1 mechanism. First, the hydroxyl group in the alcohol is protonated in a fast step to form an alkyloxonium ion. Next, a molecule of water is lost from the alkyloxonium ion in the slow, rate-determining step, leaving behind a carbocation. Finally, water, which is the conjugate base of H3O+, removes a β hydrogen from the carbocation to yield the alkene. This step regenerates the acid catalyst.

Figure2

In these reactions, the stability of the carbocation intermediate determines the major products. When possible, secondary carbocations undergo rearrangement to form more stable tertiary carbocations. Additionally, when isomeric products are possible, the more-substituted alkene, or Zaitsev's product, is favored.

Primary alcohols would yield highly unstable primary carbocations. As a result, their dehydration occurs via the E2 mechanism. This mechanism also begins with the protonation of the alcohol. In the next step, a base removes the β hydrogen, and a water molecule is lost. Thus, a double bond is formed, yielding a terminal alkene.

Figure3

However, in the acidic solution, rehydration of the double bond (according to Markovnikov's rule) followed by a 1,2-hydride shift can yield a secondary carbocation. Loss of a proton, in accordance with Zaitsev's rule, then results in a mixture of the terminal and rearranged alkenes.

Figure4

Secondary or tertiary alcohols can also undergo dehydration via the E2 mechanism if the hydroxyl group is first converted to a good leaving group, such as a tosylate. Treatment of the tosylate with a strong base yields the alkene.

 Core: Organic Chemistry

Overview of Skeletal Muscle

JoVE 12521

Skeletal muscles are composed of a bundle of muscle fibers and are attached to bones through tendons. Each skeletal muscle fiber is a single muscle cell. The sarcolemma, the plasma membrane of a skeletal muscle cell, consists of a lipid bilayer and glycocalyx that supports muscle fibers. The sarcolemma extends into the muscle cells to form tubular structures called transverse or T-tubules. Each side of the T-tubules consists of a membrane-bound structure called the sarcoplasmic reticulum, similar to the smooth endoplasmic reticulum. The storage and release of calcium ions in response to a neuronal signal.

Contractile elements called sarcomeres aid the contraction and relaxation of skeletal muscle fibers to control various body movements, protect internal organs, and maintain body posture. A sarcomere mainly consists of actin and myosin protein filaments. In a relaxed muscle cell, the active binding site of actin is masked by a protein called tropomyosin. This masking prevents actin from binding to myosin, preventing muscle contraction. When a nerve impulse stimulates a muscle fiber, the sarcoplasmic reticulum releases calcium into the sarcoplasm. The released calcium binds to a protein called troponin on the thin filaments of muscle fiber. This induces structural changes in troponin, allowing it to attach to tropomyosin. The binding of troponin to tropomyosin relieves the inhibition of actin and enables myosin to bind to actin, further causing muscle contraction.

 Core: Cell Biology

Plant Cell Wall

JoVE 13365

Plant cells have a cell wall, a rigid outer covering that protects the cell and provides shape and support. During cell division, a mixture of enzymes, proteins, and glucose molecules is transported via vesicles to the center of the cell. These vesicles continuously fuse and build a cell plate between the dividing cells. As the cell plate matures, new polysaccharides are added to it to form the cell walls of the daughter cells. The predominant polysaccharide in the cell wall is cellulose, made up of repeating glucose units. As a cell matures, its cell wall specializes according to its cell type. For example, the parenchyma cells of leaves possess only a thin, primary cell wall. Collenchyma and sclerenchyma cells, on the other hand, mainly occur in the outer layers of a plant's stems and leaves. These cells give the strength and support by either partially thickening their primary cell wall (i.e., collenchyma) or depositing a secondary cell wall (i.e., sclerenchyma). The varying cell wall compositions determine the function of specific cells and tissues.

Some plants, such as trees and grasses, deposit a secondary cell wall around mature cells. Secondary cell walls typically contain three distinct layers. In each layer, the cellulose microfibrils are organized in different orientations.

Transportation through pits

All plant cell walls have small holes, or pits, that allow water transport, nutrients, and other molecules. In a pit, the middle lamella and primary cell wall merely form a thin membrane that separates adjacent cells. The secondary cell wall may be deposited around the pit but not within.

Rigidity against turgor pressure

Plant cells absorb water and nutrients and store them in the vacuole. As the vacuole expands, it pushes the plasma membrane against the cell wall. This turgor pressure supports the upright and rigid structure of plants. The cell wall prevents the cells from rupturing under this pressure.

Storage

In addition to providing structure and support, plant cell walls also store nutrients. Seeds, for example, may store sugars in the cell walls of cotyledon and endosperm tissues for use during early plant growth. The cell wall is the principal barrier and defense against pathogenic bacteria, viruses, and fungi.

Adapted from Section 4.3 Eukaryotic cells and 30.1 Plant biology Openstax biology 2e

 Core: Cell Biology

Two-dimensional Gel Electrophoresis

JoVE 13382

Two-dimensional gel electrophoresis is a high-resolution protein separation method first introduced by O' Farrell and Klose in 1975. This method involves protein separation by two dimensions, mass and charge, making it more accurate than one-dimensional gel electrophoresis.

The first dimension separation uses the isoelectric focusing or IEF technique performed on immobilized pH gradient (IPG) strips that separate proteins according to their isoelectric points.

Biological samples, such as  cells and tissues, are first prepared for IEF by treatment with a buffer containing urea, dithiothreitol, detergents, and ampholytes. Here, urea and detergent solubilize and denature the proteins. Dithiothreitol is a reducing agent that cleaves disulfide bonds. Ampholytes used in a suitable concentration solubilize the protein and maintain pH. The buffer also rehydrates the IPG strip that helps it absorb the proteins before their charge-based separation.

Proteins are amphoteric molecules with positive and negative charges. The amphoteric protein molecules migrate on the pH gradient of the IPG strip under the electric current. Once the proteins reach their isoelectric point, they immobilize because, at this point, they have no net charge. IPG strips have a fixed pH gradient; for instance, strips with a smaller pH range (for example, pH 4-7) can detect proteins with higher accuracy than those with a broader range (pH 3-10). Once the proteins separate, the IPG strips are treated with sodium dodecyl sulfate (SDS) buffer to proceed to the next phase, i.e., mass-based separation using SDS-PAGE (polyacrylamide gel electrophoresis). SDS covers the uncharged protein on the IPG strips with a negative charge, and glycerol reduces electroendosmosis which helps the transfer of proteins from the first to the second dimension.

Generally, bands on the gel can be visualized upon staining. Sometimes radiolabelled samples are used to obtain an autoradiographic image. Samples capable of producing fluorescent images can also be used.

Two-dimensional gel electrophoresis has gained importance owing to the precision of the outcome results. However, the method is unsuitable for distinguishing highly hydrophobic proteins,  proteins with high molecular mass, or those present in very low quantities per cell.

 Core: Cell Biology

Super-resolution Fluorescence Microscopy

JoVE 13398

Super-resolution fluorescence microscopy (SRFM) provides a better resolution than conventional fluorescence microscopy by reducing the point spread function (PSF). PSF is the light intensity distribution from a point that causes it to appear blurred. Due to PSF, each fluorescing point appears bigger than its actual size, and it is the PSF interference of nearby fluorophores that causes the blurred image. Various approaches to achieving higher resolution through SRFM have recently been developed.

Photoactivated Localization Microscopy

Photoactivated localization microscopy (PALM) is a type of fluorescence microscopy that captures high-resolution images using a single-molecule detection and localization approach. For example, two fluorescent spots 75 nm apart may appear to be a single spot during imaging due to interference of their PSFs. In such cases, especially for live-cell imaging, PALM is a suitable technique to resolve the interference and provide a better resolution.

In PALM, a variant of green fluorescent protein (GFP) with a different excitation wavelength and high fluorescence is employed. In the first step, a few GFPs are activated and imaged with very high precision. In the next step, another set of GFPs is activated and imaged. Step by step, all the GFPs across the specimen are thus recorded. Finally, the data is processed to generate a high-resolution image.

Stochastic Optical Reconstruction Microscopy

In stochastic optical reconstruction microscopy (STORM), unique photo-switchable probes are used for specimen imaging. The emission from these probes can be switched on and off using lights of different wavelengths. Thus, only a few fluorophores can be activated at a point in time such that the number of activated fluorophores is significantly lesser than the number of deactivated fluorophores. This selective activation of probes helps determine their precise position in the specimen. Following this, the center of each fluorescent probe is identified and marked. The process is then repeated to record all the fluorescent probes in the specimen. Finally, a high-resolution composite image is constructed by superimposing these multiple images.

 Core: Cell Biology

Source And Potency Of Stem Cells

JoVE 13462

Stem cells are undifferentiated cells with extensive self-renewal properties that help them maintain their population during the fetal and adult stages of life. They can specialize in all cell types of the human body. However, their differential potential may vary and can be classified into five types. Stem cells can be (1) Totipotent, (2) Pluripotent, (3) Multipotent, (4) Oligopotent, and (5) Unipotent. Each stem cell has a specific origin; the fertilized egg or zygote is a totipotent cell and can develop into the embryonic and extraembryonic tissues. As development progresses, the zygote undergoes repeated rounds of cell divisions and eventually loses its totipotency to become a 64- cell mass called a blastocyst. The blastocyst consists of the rapidly proliferating inner cell mass (ICM) and the trophectoderm (TE) outer layer. Both ICM and TE are sources of pluripotent stem cells. ICM can differentiate into any embryonic tissue, including blood cells, muscle cells, neurons, skin, and intestinal cells, but not the extraembryonic tissues. In contrast, TE cells can only form the placenta, amnion, and chorion.   ICM is also called embryonic stem cells or ES cells. ES cells are routinely used in regenerative medicine to treat stroke and neurodegenerative diseases such as Parkinson's or Alzheimer's.

Multipotent stem cells include adult stem cells or ASCs. They reside in the specific tissue type and produce cells of that particular tissue lineage only. For example, ASCs of the bone marrow are called hematopoietic stem cells or HSCs, which can differentiate into all blood cells but not gut cells or muscle cells. Oligopotent stem cells have further restricted tissue lineage and can differentiate to only fewer cell types of the tissue. The bronchoalveolar stem cells (BASCs) of the lungs are oligopotent and can differentiate into the bronchiolar epithelium or alveolar epithelium. Unipotent cells can only generate one kind of cell, such as muscle stem cells. They mature into only muscle cells in the body.

 Core: Cell Biology

Induced Pluripotent Stem Cells

JoVE 13479

Stem cells are undifferentiated cells that divide and produce different cell types. Ordinarily, cells that have differentiated into a specific cell type are terminally differentiated; however, scientists have found a way to reprogram these mature cells so that they dedifferentiate and return to an unspecialized, proliferative state. These cells are pluripotent like embryonic stem cells—able to produce all cell types—and are called induced pluripotent stem cells (iPSCs).

Somatic cells are reprogrammed into iPSCs using viral vectors to deliver the genes encoding Oct4, Sox2, Klf4, and c-Myc (OSKM) transcription factors. There are various ways to confirm the transformation. Among these, checking for teratoma formation is a routinely performed test.  A teratoma is an encapsulated, benign tumor formed by the successfully converted iPSCs, consisting of tissues from all three germ layers– endoderm, mesoderm, and ectoderm. Other characteristics of iPSCs include unlimited proliferation, long telomeres, and expression of pluripotent stem cell markers such as Oct4, c-Myc, and Nanog.

iPSCs are potentially valuable in medicine because a patient could receive a transplant of the required cells generated from another cell type from their own body. This is called autologous transplantation, and it reduces the risk of transplant rejection that can occur when tissues are transplanted between individuals.

Early Clinical Trials

The first clinical trial using iPSCs involved transplanting retinal cells derived from iPSCs into patients with age-related macular degeneration–a disease of the retina affecting normal vision. Since then, several iPSC clinical trials have been approved to treat Parkinson’s disease, heart disease, and spinal cord injury. Cells taken from patients and turned into iPSCs are also used to study their diseases in the laboratory. In general, iPSCs provide another source of stem cells for scientific research.

Limitations and Risks

Generating iPSCs from somatic cells is a long and inefficient process, which may take up to 4 weeks, with only 0.01 to 1% of somatic cells being able to form pluripotent stem cells. Secondly, all the OSKM factors must be expressed in a given cell in optimal amounts to become a  pluripotent one.

The viral vectors for delivering the OSKM factors can introduce the genes in an undesired location in the genome, activating genes that can cause cancer. Similarly, c-Myc is reported to be an oncogene, which, if overexpressed, can also cause cancer. 

 Core: Cell Biology

Inductors

JoVE 13800

An inductor, also known as a choke, is a circuit component created to have a specific inductance. Inductors are among the crucial circuit components used in modern electronics, along with resistors and capacitors. They serve as a barrier against changes in a circuit's current. An inductor tends to suppress current changes in an alternating-current circuit that are faster than desired. In a direct-current circuit, an inductor aids in preserving a constant current despite changes in the applied emf.

Although the electric field associated with the magnetic induction effect is non-conservative, there is a real potential difference between the inductor's terminals, caused by conservative electrostatic forces.

One common application of inductance is to allow traffic signals to sense when vehicles are waiting at a street intersection. An electrical circuit with an inductor is placed in the road underneath the location where a waiting car will stop. The car's body increases the inductance, and the circuit changes, sending a signal to the traffic lights to change colors. Similarly, metal detectors used for airport security employ the same technique. A coil or inductor in the metal detector frame acts as a transmitter and a receiver. The pulsed signal from the transmitter coil induces a signal in the receiver. Any metal object in the path affects the circuit's self-inductance. Metal detectors can be adjusted for sensitivity and can also sense the presence of metal on a person.

Inductors are also crucial in fluorescent light fixtures. In such fixtures, the gas that fills the tube is ionized and glows due to the current flowing from the wiring into the gas. The higher the current, the more strongly ionized the plasma and the lower its resistance. If a sufficiently high voltage is given to the plasma, the circuitry outside the fluorescent tube can be damaged. To prevent this, an inductor or magnetic ballast is connected in series with the fluorescent bulb, which keeps the current from exceeding its limits.

 Core: Physics

Propagation of Uncertainty from Systematic Error

JoVE 14512

The atomic mass of an element varies due to the relative ratio of its isotopes. A sample's relative proportion of oxygen isotopes influences its average atomic mass. For instance, if we were to measure the atomic mass of oxygen from a sample, the mass would be a weighted average of the isotopic masses of oxygen in that sample. Since a single sample is not likely to perfectly reflect the true atomic mass of oxygen for all the molecules of oxygen on Earth, the mass we obtain from this particular sample will give rise to sampling error or a type of systematic error.This means that there is an uncertainty in the atomic mass of oxygen. It is not a random error but specific to the system or a systematic error. The probability of finding a particular value is the same within the uncertainty range. The distribution plot of values has a rectangular shape. With a sharp cut in the frequency of occurrence at the minimum and maximum reported values, the values in between have roughly the same frequency of appearance. This differs from the Gaussian distribution of values with random errors, where the frequency of occurrence is high near the central or mean value and tapers down as the value deviates outwards from the mean.

The standard deviation or standard uncertainty obtained from the rectangular distribution plot can be used to measure uncertainty. The individual uncertainties from constituent atoms in a molecule propagate to the uncertainty of its molecular mass. The propagation of systematic errors does not follow the same rules as that of random errors. Because systematic errors are not reduced by repeated measurements, their propagation is much more additive. Consider the following example: In a molecule containing multiple identical atoms, the uncertainty in the mass of one atom is multiplied by the number of similar atoms to get the total contribution of uncertainty toward the molecular mass of those particular atoms. If heteroatoms are present in a molecule, the uncertainty in the molecular mass is the square root of the sum of the squares of total uncertainties from each type of atom present.

 Core: Analytical Chemistry

Factors Affecting Activity Coefficient

JoVE 14528

The extended Debye-Hückel equation indicates that the activity coefficient of an ion in an aqueous solution at 25°C depends on three partially interdependent properties: the ionic strength of the solution, the charge of the ion, and the ion size. 

The activity coefficient value for an ion is close to one when the solution has almost zero ionic strength, i.e., when the solution shows close to ideal behavior. As the ionic strength of the solution increases from 0 to 0.1 mol/L, a decrease in the activity coefficient value is observed.

Solutions with an ionic strength above 0.1 mol/L are not well represented by the Debye-Hückel equation. In such solutions, the activity coefficient of an ion may exceed unity. The ionic strength of the solution is closely related to the charge of the ions present in the solution, thereby affecting their activity coefficient. For a particular ionic strength, ions with less charge have activity coefficient values higher than ions with a higher charge. In other words, the deviation from ideality is more pronounced for multiply charged ions than for singly charged ions.

Ions with similar charges often have similar activity coefficients. The dominant factor in the coefficient discrepancies here is the ion size parameter. Ion size describes the effective diameter of a hydrated ion, and even ions with similar charges tend to have different extents of hydration and therefore different ion sizes. Ions with similar charges can be hydrated to different extents, resulting in different effective diameters. These diameters are known as the ion size parameters. Ions with smaller ion size parameters deviate more from ideality, resulting in a lower activity coefficient than those with larger ion size parameters but the same charge.

 Core: Analytical Chemistry

Titration in Nonaqueous Solvents

JoVE 14544

Most acid-base titrations are performed in an aqueous medium. In aqueous titrations, water competes with weaker acids or bases for proton donation or acceptance, leading to ambiguous endpoints in the titration curve. Water also affects the partial ionization of weak acids or bases. For example, water accepts a proton from acetic acid to form hydronium and acetate ions. The hydronium ion formed is a stronger acid than acetic acid, and the acetate ion is a stronger base than water. As a result, they react to give back the reactants. The effect of this process on weak acids and bases means that calculations based purely on pKa values may yield inaccurate results.

Using a non-aqueous solvent like ammonia, which is a stronger base than water, enables the complete ionization of acetic acid into acetate ion, effectively turning acetic acid into a strong acid in ammonia. Put another way, ammonia has a higher dissociation constant (Ks) than water, thereby increasing the equivalence point and sharpening the endpoint in the titration curve of acetic acid. This is commonly observed in the titration of weak acids and bases in non-aqueous solvents. Note that Ks of water is denoted by Kw.

There are four types of non-aqueous solvents: aprotic–unable to donate protons, protophilic–able to accept protons, protogenic–able to donate protons, and amphoteric–able to donate and accept protons. In addition to their uses in the titration of weak acids or bases, some of these solvents can also be used effectively in titrating organic analytes, which have poor solubility in water. The reactions occurring in non-aqueous titrations are explained by the Bronsted-Lowry theory of acids and bases. Here, while an acid behaves as a proton donor, the base is a proton acceptor.

 Core: Analytical Chemistry

Precipitation Gravimetry

JoVE 14583

Precipitation gravimetry is based on converting an analyte into a sparingly soluble precipitate, which is separated by filtration and weighed. An ideal precipitate should be pure, insoluble, of known composition, and easily filtered from the reaction mixture.

In determining nickel by gravimetric analysis, a precipitant of ethanolic dimethylglyoxime is added to a hot nickel salt solution. This is quickly followed by the dropwise addition of dilute ammonia solution until precipitation occurs. A slight excess of ammonia is added to ensure complete precipitation, followed by digestion for a few hours over a hot water bath to remove any impurities that may have occluded into the precipitate by dissolving and slowly reforming it multiple times. The resultant solution is allowed to stand for some time and cooled. The cold solution is filtered using a sintered glass funnel to obtain the red precipitate. This precipitate is washed with cold water, and the sintered glass funnel is dried in a hot air oven for an hour. The dried precipitate is cooled in a desiccator and weighed. From the mass of nickel dimethylglyoximate, the nickel content is calculated using the stoichiometric mole relationship.

 Core: Analytical Chemistry

Centrifugation

JoVE 14616

Centrifugation is a separation technique based on differences in density or size. It is commonly used to separate solids from aqueous interferents. During centrifugation, the sample is placed in centrifugation tubes and spun at high angular velocity, which allows centrifugal force to act differentially on the different densities or masses of the components. After spinning, the supernatant liquid is decanted. Depending on the specific application, either the pellet or the supernatant is retained for further purification. Sometimes, the supernatant is subjected to additional rounds of centrifugation and downstream processing.

In principle, centrifugation separates particles based on differences in size and density. Larger particles tend to be heavier and tend to sediment first. If particles are similar in size, denser particles tend to sediment first due to higher sedimentation rates. In practice, complex mixtures such as cell lysates may not yield distinctly separated particles despite differences in density or size. To address this issue, researchers have devised a variety of centrifugation techniques.

The simplest centrifugation technique is differential centrifugation. Here, the particles to be separated have similar densities, so larger particles will sediment at lower speeds. The speed is increased stepwise until the target particles are isolated.

A more sophisticated method is equilibrium density gradient centrifugation. In this technique, the analyte is placed in a solution with a preformed density gradient or a solution that forms a density gradient during centrifugation. Here, the density of the solution increases towards the bottom of the tube, so the sedimentation rate of each analyte component decreases as it moves toward the bottom. When a component's density equals that of the solution, the centrifugal force acting on it becomes zero, and sedimentation ceases. As a result, each component is isolated in a layer equal to its density.

 Core: Analytical Chemistry

Relaxation of Skeletal Muscles

JoVE 14843

The period of muscle contraction primarily influences the duration of stimulation at the neuromuscular junction (NMJ), the presence of free calcium ions in the sarcoplasm, and the availability of energy or ATP to support contractions.

When an action potential reaches the axon terminal, it depolarizes the membrane and opens voltage-gated sodium channels. Sodium ions enter the cell, further depolarizing the presynaptic membrane. This depolarization causes voltage-gated calcium channels to open. Calcium ions entering the cell initiate a signaling cascade that causes small membrane-bound vesicles, called synaptic vesicles, containing neurotransmitter molecules to fuse with the presynaptic membrane. The fusion of a vesicle with the presynaptic membrane causes neurotransmitters to be released into the synaptic cleft, the extracellular space between the presynaptic and postsynaptic membranes.

A single action potential reaching the presynaptic end of the neuromuscular junction triggers the release of acetylcholine in the synaptic cleft. Once the receptors present on the postsynaptic membrane recognize acetylcholine, the neurotransmitter quickly dissociates from the receptor and is degraded by acetylcholinesterase present in the synaptic cleft. The swift removal and inactivation of acetylcholine from the receptor prevent overstimulation of the postsynaptic membrane from excessive firing of action potentials. As a result, the muscle membrane repolarizes, and the calcium concentration in the muscle fiber returns to a normal resting level. There are two mechanisms involved in restoring normal intracellular calcium levels: the active transport of calcium into the sarcoplasmic reticulum (SR) and of calcium across the sarcolemma into the extracellular fluid. Of the two mechanisms, the transport into the SR is more important. The SR promptly resumes normal permeability and actively reabsorbs calcium ions from the surrounding cytosol. When calcium concentration in the cytosol decreases, the contraction process reverses. Troponin releases calcium ions and reverts to its original conformation. As a result, the tropomyosin moves back, covering the myosin-binding sites on actin. Therefore, a single action potential only has a brief impact on the muscle fiber.

When there is a rapid succession of action potentials at the presynaptic end of the neuromuscular junction, there is a continuous release of acetylcholine. This, in turn, triggers a series of action potentials in the sarcolemma at the post-synaptic end and results in continuous high levels of calcium in the cytosol. Under these circumstances, the contraction cycle repeats itself without muscle fibers undergoing relaxation.

 Core: Anatomy and Physiology

Axial and Appendicular Muscles

JoVE 14865

Skeletal muscles, the key players in our body's movement, can be classified into two groups based on their location and function: axial muscles and appendicular muscles. These classifications reflect the primary roles the muscles play in the body's structure and movement.

Axial Muscles

Axial muscles, situated along the body's midline, are intricately connected to the axial skeleton, which includes the skull, spine, ribs, and sternum. These muscles facilitate facial expressions and play a crucial role in various head and neck movements, speaking, and eating. They even contribute to controlling the movements of the vertebral column in coordination with the back muscles. Furthermore, the thoracic muscles, including the rectus, oblique, and transverse abdominal muscles, form the anterolateral walls of the trunk, offering stability and support. Additionally, the lower axial skeleton muscles extend between the sacrum and pelvic girdle, actively participating in forming the pelvic floor.

Appendicular Muscles

In contrast, appendicular muscles are predominantly associated with the pectoral girdle, pelvic girdle, and limbs. These muscles provide stability and support and actively contribute to the body's movement during activities like walking, running, and various physical actions. The upper limbs are equipped with muscles that extend from the shoulders down to the hands, enabling us to perform intricate tasks. Similarly, the lower limbs have muscles stretching from the hips to the feet, ensuring mobility and balance.

 Core: Anatomy and Physiology

Organization of the Nervous System

JoVE 14882

The nervous system is one of the most complex systems in our body. It is organized into two main divisions: the central nervous system (CNS) and the peripheral nervous system (PNS).

The CNS, comprising the brain and spinal cord, houses billions of neurons. The brain is housed in the skull, while the spinal cord is linked to the brain through the foramen magnum of the occipital bone and is surrounded by the protective structure of the vertebral column. It is responsible for processing various types of sensory information, as well as being the source of thoughts, emotions, and memories. It also initiates signals that trigger muscle contractions and gland secretions.

On the other hand, the PNS includes all nervous tissue outside the CNS. It comprises nerves and sensory receptors—structures that monitor changes in both external and internal environments.

The PNS is further divided into sensory and motor divisions. The sensory division brings input to the CNS from sensory receptors throughout the body, providing sensory information about both somatic and special senses. The motor division, which transmits output from the CNS to effectors like muscles and glands, is divided into the somatic nervous system (SNS) and the autonomic nervous system (ANS). The SNS controls the activity of skeletal muscles, while the ANS controls the activities of smooth muscles, cardiac muscles, and glands.

 Core: Anatomy and Physiology

Integration of Synaptic Events

JoVE 14899

Synaptic integration mainly includes the summation of graded potentials. Graded potentials, regardless of their type, cause subtle alterations in membrane voltage, resulting in either depolarization or hyperpolarization. These incremental changes, when combined or summed, can propel the neuron toward its threshold. Consider, for example, a membrane experiencing a +15 mV shift, causing it to depolarize from -70 mV to -55 mV. In this scenario, graded potentials govern the membrane's ability to reach the threshold.

Synaptic Integration

The process by which multiple synaptic potentials combine within one postsynaptic neuron is known as synaptic integration. It is crucial for neural computation, as it influences how neurons process and transmit information. There are three main types of synaptic integration:

  1. Temporal Summation: Temporal summation occurs when a single presynaptic neuron fires many times in succession, causing the postsynaptic neuron to reach its threshold and fire an action potential. This type is dependent on the timing of the incoming signals.
  2. Spatial Summation: Spatial summation involves inputs from multiple presynaptic neurons firing at the same time to cause the postsynaptic neuron to fire. This type is dependent on the physical location of the incoming signals.
  3. Dendritic Integration: Dendrites, the branches off of neurons, also play a role in synaptic integration. Dendritic integration refers to the combination of inputs arriving at different parts of the dendrite. This form of integration is complex due to the active properties of dendrites, which can amplify or attenuate the incoming signals.

Significance for Neural Computation and Neural Outputs

Synaptic integration plays a significant role in neural computation as it helps to determine whether a neuron will fire an action potential or not. By summing the excitatory and inhibitory inputs, the neuron can make a 'decision' about whether to pass the signal along. The type of integration (temporal, spatial, or dendritic) can result in different types of neural outputs, affecting everything from sensory perception to motor control. Thus, understanding synaptic integration is key to understanding how the brain processes and responds to information.

The integration of synaptic events plays a crucial role in the development and functionality of the nervous system. It facilitates the transfer, processing, and transmission of information between neurons, allowing for effective communication within neural networks and coordination of complex behaviors. Moreover, synaptic integration enables intercommunication among different regions of the brain. For instance, when individuals learn new skills, various brain regions exchange information to aid in task completion. This interplay between brain regions is pivotal for memory formation, decision-making, and behavior.

Integration of synaptic events also plays a role in modulating neural activity patterns in response to external stimuli. This helps the brain detect changes in the environment and respond appropriately. For example, when a person hears an alarm clock, their brain can recognize this stimulus as a cue to wake up and initiate the appropriate behaviors. Synaptic integration is also essential for homeostasis; it enables neurons to adjust their activities based on incoming signals from other neurons. This helps the body maintain a state of equilibrium and protects it from external influences that could disrupt its normal functioning.

 Core: Anatomy and Physiology

Functional Brain Systems: Limbic System

JoVE 14915

The limbic system, often called the "emotional brain," is a complex set of structures located deep within the brain. The intricate network of the limbic system supports a wide range of psychological functions, from emotional regulation to memory formation and sensory processing. This functional brain region encompasses specific parts of the diencephalon and the cerebrum, integrating the higher mental functions of the cerebral cortex with the primitive emotional responses of the deep brain structures.

The Gyri

The limbic lobe includes three prominent gyri — the cingulate gyrus, the parahippocampal gyrus, and the subcallosal gyrus.

  • • Cingulate Gyrus: The cingulate gyrus is located in the medial aspect of the cerebral hemisphere, positioned above the corpus callosum. It is involved in processing emotions and regulating behavior. It acts as a pathway connecting the amygdala and the hippocampus, which is crucial for emotional response and memory formation.
  • • Parahippocampal Gyrus: This gyrus is situated in the inferior region of the temporal lobe next to the hippocampus, serving as an essential interface between the limbic system and sensory inputs. It comprises several subregions, notably the entorhinal, perirhinal, and postrhinal cortices, each contributing unique functions to memory processing and spatial navigation. Its extensive connections with the hippocampus establish the entorhinal-hippocampal circuit, a fundamental memory encoding and retrieval pathway. Through this intricate network, the parahippocampal gyrus facilitates the complex processes of spatial navigation and recognition memory.
  • • Dentate Gyrus: The dentate gyrus is located between the parahippocampal gyrus and the hippocampus. It forms part of the hippocampal formation, primarily associated with memory formation and spatial navigation. It is characterized by its densely packed granule cells and generates new neurons throughout life via adult hippocampal neurogenesis. These newly formed neurons are believed to contribute to learning and memory processes, particularly pattern separation, which is the ability to distinguish between similar experiences or stimuli.

The Amygdala

Another vital component, the amygdala, is an almond-shaped nucleus within the temporal lobe known for its role in processing emotions, especially fear and anxiety. It is involved in detecting and evaluating threats in the environment, as well as in emotional learning and memory formation. The amygdala has extensive connections with other brain regions, including the prefrontal cortex, hippocampus, and hypothalamus. These connections allow the amygdala to integrate sensory information and emotional cues from the environment and coordinate appropriate behavioral and physiological responses. Its functions are critical for survival, as it generates emotional responses to threats and helps store memories of events that triggered strong emotions.

Additional Structures

In addition to these cortical areas, the limbic system includes several other structures contributing to its complex functionality. The mammillary bodies within the hypothalamus are integral for memory recall and spatial memory, indicating their importance in navigating and remembering spaces. Furthermore, the anterior and medial nuclei of the thalamus serve as vital relay stations, channeling information into the cerebral cortex and supporting emotional regulation and memory processes. The olfactory bulbs play a pivotal role in the sensory processing of smells, connecting specific scents to memories and emotions, thereby demonstrating the limbic system's significant involvement in our emotional responses to different odors.

 Core: Anatomy and Physiology

Sensory Perception: Organization of the Somatosensory System

JoVE 14933

The somatosensory system is the central and peripheral nervous system component that senses and processes touch, pressure, pain, temperature, and body position or proprioception. The process of sensation takes place at three levels:

The receptor level:

The receptor level is the first stage of sensation. It involves the detection of a stimulus by specialized sensory receptors. The stimulus must arrive within the receptor's receptive field. Next, the receptor converts the energy of the stimulus into an electrical signal via transduction, which leads to the generation of membrane potential in the receptor cell. When the potential reaches a certain threshold, a nerve impulse is generated

The circuit level:

The circuit level is the second stage of sensation. The generated nerve impulse now travels to the central nervous system (CNS). Several ascending tracts carry the impulses from the receptors to their final destination in the CNS. The stimulus from the facial area is carried and transmitted by cranial nerves, while stimulus from the back of the head and rest of the body travels by spinal nerve. Sensory neurons conducting impulses from the peripheral nervous system (PNS) to the CNS are called first-order neurons.

The perceptual level:

The perceptual level is the third and final stage of sensation. It involves the reception of sensory information by different regions of the CNS, depending on the stimuli. Only the impulses processed in the cerebral cortex are consciously perceived.

 Core: Anatomy and Physiology

Sympathetic Pathways: Collateral Ganglia and Adrenal Medulla

JoVE 14950

The sympathetic pathways of the collateral ganglia and adrenal medulla serve unique but interconnected roles in the sympathetic response.

Collateral Ganglia

Sympathetic preganglionic axons reach the collateral ganglia along the route of splanchnic nerves. These nerves bypass the sympathetic trunk and communicate with sympathetic postganglionic neurons housed in the prevertebral ganglia. These ganglia supply the organs of the abdominopelvic cavity.

The greater splanchnic nerve, formed by the preganglionic axons from the thoracic ganglia (T5–T9 or T10), interfaces with the stomach, spleen, liver, kidneys, and small intestine. The lesser splanchnic nerve (T10–T11) reaches the aorticorenal and the superior mesenteric ganglia, affecting the small intestine and proximal colon. The lumbar splanchnic nerve, formed by axons from the lumbar ganglia (L1–L4), interfaces with the distal colon, rectum, urinary bladder, and genital organs through the inferior mesenteric ganglion.

Adrenal Medulla

The adrenal medulla pathway, on the other hand, is unique. Some sympathetic preganglionic axons traverse the sympathetic trunk, greater splanchnic nerves, and the celiac ganglion to reach the adrenal medullae without synapsing. This area houses chromaffin cells, which are modified sympathetic postganglionic neurons. But, unlike typical neurons, they lack dendrites and axons, and instead of transmitting to another organ, these cells discharge hormones into the bloodstream. Upon stimulation by sympathetic preganglionic neurons, chromaffin cells release catecholamine hormones, predominantly epinephrine (80%) and norepinephrine (20%), with dopamine in trace amounts. These hormones disperse into the bloodstream, intensifying the body's responses initiated by sympathetic postganglionic neurons.

 Core: Anatomy and Physiology

Anatomy of the Ear

JoVE 14969

Auditory sensation, commonly called hearing, involves the transformation of sonic waves into neural impulses facilitated by the structures of the auditory organ. The prominent, flesh-like structure on the side of the head, called the auricle, directs sound waves towards the auditory canal. The auricle is often mislabeled as the pinna, a term more aligned with mobile structures like a feline's external ear. The auditory canal penetrates the cranium via the external auditory meatus of the temporal bone and culminates at the tympanic membrane. The tympanic membrane, more colloquially known as the eardrum, vibrates when impacted by sound waves. Collectively, the auricle, ear canal, and tympanic membrane form the external ear.

The middle ear is comprised of three small ossicles, or bones. These structures are called the malleus, incus, and stapes, derived from Latin and translated to mean hammer, anvil, and stirrup, respectively. The malleus, connected to the eardrum, meets with the incus, which connects to the stapes. The stapes link with the inner ear, where sound waves are transformed into neural signals, a pivotal point in the auditory process. The middle ear communicates with the pharynx via the Eustachian tube, which balances air pressure on either side of the tympanic membrane. This tube is generally closed, opening only when the pharyngeal muscles contract during swallowing or yawning.

The inner ear, characterized by a labyrinthine structure owing to a series of canals within the temporal bone, is subdivided into two sections. The two subsections are the cochlea and the vestibule, which facilitate hearing and balance, respectively. The neural impulses from these regions are relayed to the brainstem via separate fiber bundles from the inner ear to the brainstem as the vestibulocochlear nerve. Sound transformation into neural signals occurs within the inner ear's cochlear region, which houses the spiral ganglia's sensory neurons. The ganglia within the spiral-shaped cochlea of the inner ear is affixed to the stapes via the oval window.

The oval window forms the commencement of a fluid-filled conduit within the cochlea termed the scala vestibuli. Extending from the oval window, the scala vestibuli traverses above the cochlear duct, the median cavity of the cochlea that hosts the auditory-transducing neurons. The scala vestibuli envelop the cochlear duct near the tip of the cochlea. The fluid-filled conduit returning to the base of the cochlea is known as the scala tympani. Beneath the cochlear duct, the scala tympani ends at the round window, sealed by a membrane that encloses the fluid within the scala. The vibrations of the ossicles, transmitted through the oval window, cause the fluid within the scala vestibuli and scala tympani to undulate. The fluid waves' frequency corresponds with the sound waves' frequency. The membrane sealing the round window protrudes or invaginates in response to the fluid motion within the scala tympani.

 Core: Anatomy and Physiology

The Parathyroid Glands

JoVE 14985

The two pairs of parathyroid glands embedded within the posterior surface of the thyroid gland are restricted by a dense capsule around them. These glands comprise two distinct cell populations—parathyroid oxyphil and parathyroid principal cells- pivotal in calcium homeostasis.

Oxyphil cells, whose functions remain elusive, emerge during late puberty, adding a layer of complexity to the parathyroid gland's intricacies. In contrast, principal parathyroid cells undertake a vital role by producing the parathyroid hormone (PTH), also known as the parathormone. These cells actively monitor circulating calcium concentrations, exhibiting a dynamic response when levels decline.

Upon detecting decreased calcium levels, principal parathyroid cells release PTH, instigating a cascade of events to mobilize calcium. PTH regulates osteoblast and osteoclast activity in the bones. It initiates the secretion of the receptor activator of a growth factor called the nuclear factor-kappa-B ligament (RANKL) by osteoblasts. RANKL, in turn, amplifies osteoclast number and activity, intensifying mineral turnover and calcium release through bone matrix erosion.

Simultaneously, the heightened blood calcium levels trigger increased kidney reuptake and stimulate calcitriol secretion, reinforcing PTH activity. This feedback loop ensures a finely tuned regulation of calcium levels within the body, exemplifying the parathyroid glands' indispensable role in maintaining systemic mineral homeostasis.

 Core: Anatomy and Physiology

Nodal Analysis

JoVE 15047

Nodal analysis is a fundamental method in electrical engineering used to simplify the process of circuit analysis. This method revolves around the concept of using node voltages as the primary variables for circuit analysis. The objective is to determine the voltage at each node in a circuit, which can then be used to find other quantities of interest, such as currents through specific components.

Consider, for instance, a simple circuit composed of three nodes and three resistors, as shown in Figure 1. The first step in nodal analysis involves selecting a reference or datum node. This node is typically chosen based on convenience, and its voltage is assigned a value of zero.

Figure1

Figure 1

The subsequent nodes, referred to as non-reference nodes, are assigned nodal voltages relative to this reference node. In the example considered here, there are two non-reference nodes labeled 1 and 2, each with their respective node voltages.

To establish a relation between the branch currents, Kirchhoff's Current Law (KCL) is applied to the non-reference nodes. KCL states that the algebraic sum of currents entering a node (or a closed boundary) is zero. This law is grounded on the principle of charge conservation – that is, a charge cannot be created or destroyed.

Following the application of KCL, Ohm's Law is used to express the branch currents passing through the three resistors in terms of the node voltages. Ohm's Law postulates that the current passing through a conductor between two points is directly proportional to the voltage across the two points.

With the branch currents expressed in terms of node voltages, these values are substituted into the equations derived from KCL. This substitution results in two simultaneous equations since, for a circuit with 'n' nodes, 'n-1' independent equations are obtained.

Lastly, if the values of the resistors and source currents are known, they can be substituted into the two equations. Solving these equations will yield the node voltages. This information is invaluable, as it can help in understanding the behavior of the circuit and in designing or troubleshooting electrical circuits.

In conclusion, nodal analysis is a powerful tool in circuit analysis, providing a systematic method to determine the distribution of voltages within a circuit, which can then be used to calculate other parameters like currents and power.

 Core: Electrical Engineering

Kirchhoff's Current Law

JoVE 15070

In the realm of electrical engineering, physicist Gustav Robert Kirchhoff made a significant contribution in 1847 by introducing Kirchhoff's laws for electric circuit analysis. These laws, particularly Kirchhoff's Current Law (KCL), have become foundational principles in understanding and analyzing electrical circuits.

Kirchhoff's Current Law is based on the principle of charge conservation. It states that at any node (a point where two or more circuit elements meet) in an electrical circuit, the total current entering the node is equal to the total current leaving the node. This means that no current is lost at the junction, reflecting the fundamental principle that electric charge is conserved.

In the application of KCL, currents leaving the node are assigned a negative sign, while currents entering the node are given a positive sign. By algebraically summing these currents at a node and integrating them over time, one can determine the total electric charge at the node.

According to the law of charge conservation, the net electric charge at any node remains constant over time, ensuring that the node stores no net charge. Therefore, if the net charge at a node is zero, the total current entering and leaving that node must also be zero. This validates the application of KCL in circuit analysis.

Interestingly, KCL can also be generalized for a closed boundary by conceptualizing a node as a closed surface condensed to a point. This allows for the application of KCL in more complex circuit environments.

KCL is also useful in determining the combined current from parallel current sources. By algebraically summing the individual currents at a node, the total current through that node can be found.

However, for KCL to hold true, a critical condition must be met: a circuit cannot have two unequal currents in series. This is because, in a series connection, the same current flows through all components, and any discrepancy would violate the law of charge conservation.

 Core: Electrical Engineering

Series and Parallel Capacitors

JoVE 15086

Capacitors, fundamental components in electronic circuits, can be connected in series and/or parallel configurations. Each configuration has different impacts on the overall behavior of the circuit.

First, consider capacitors connected in series to a battery. In this configuration, the plate connected to the battery's positive terminal develops a positive charge, while the plate attached to the negative terminal becomes negatively charged. An equal magnitude of charge is induced on the other plates, illustrating that the same current flows through each capacitor in a series connection.

Applying Kirchhoff's voltage law, which states that the total voltage around any closed loop in a circuit must equal zero, to the loop and substituting for the voltage across each capacitor yields an expression for the total voltage across the equivalent capacitor. In this case, the initial voltage across the equivalent capacitor equals the sum of the initial voltages across each capacitor. The reciprocal of the equivalent capacitance in a series connection is the sum of the reciprocals of the individual capacitances.

On the other hand, when capacitors are connected in parallel, the potential difference developed across each capacitor equals the battery voltage. This is because all the capacitors in a parallel configuration share the same voltage source.

By applying Kirchhoff's current law, which states that the total current entering a junction must equal the total current leaving it, and substituting the current for each capacitor, we can determine the total current flowing through the equivalent capacitor. In a parallel connection, the equivalent capacitance is simply the sum of the individual capacitances.

In conclusion, understanding how capacitors behave in series and parallel configurations is crucial in electronics. This knowledge allows engineers to manipulate the total capacitance of a circuit and, as a result, control the circuit's response to different signals. Whether in filtering applications, power supply smoothing, or signal coupling, capacitors and their configurations play a vital role.

 Core: Electrical Engineering

Sinusoidal Sources

JoVE 15104

Direct current (DC) refers to an electric current that flows in a single direction, maintaining a constant polarity. This is in contrast to alternating current (AC), which periodically changes its direction and magnitude. AC forms the backbone of modern electricity transmission and distribution systems due to its efficient long-distance transmission capabilities.

In homes, the power supplies use sinusoidal sources to provide electricity. These sources generate a voltage that varies sinusoidally over time. In other words, the voltage oscillates between a maximum and minimum value in a smooth, wave-like pattern.

Mathematically, this variation is represented by a harmonic function known as a sinusoid. The sinusoid is defined by three primary parameters: its amplitude (the peak value), angular frequency (the rate at which it oscillates), and phase (its position relative to a reference point).

The sinusoidal waveform repeats its pattern every T seconds, where T is referred to as the period of the function. The period is inversely related to the angular frequency, representing the duration of one complete cycle of the waveform.

On the other hand, the sinusoid frequency is defined as the number of cycles that occur per second. It is the reciprocal of the period and is typically measured in Hertz (Hz).

To understand the concept of phase, consider two sinusoids with the same frequency. If one sinusoid reaches its maxima and minima earlier than the other, then the first sinusoid is said to "lead" the second. In this case, the two sinusoids are out of phase, and the shift between them is known as the phase difference.

On the contrary, if the two sinusoids reach their maxima and minima simultaneously, they are considered to be in phase. This means they have a phase difference of zero.

Understanding these fundamental concepts related to AC is crucial for various applications, from designing electronic circuits to analyzing power systems. The principles of AC underpin many technologies in our daily lives, including household appliances, telecommunications systems, and power generation facilities.

 Core: Electrical Engineering

Design Example: Automobile Ignition System

JoVE 15173

The automobile's ignition system plays a vital role by ensuring the timely ignition of the fuel-air mixture in each cylinder. This ignition is facilitated by a spark plug, which is composed of two electrodes separated by an air gap. A spark forms across this air gap when a substantial voltage is generated between the electrodes, leading to the ignition of the fuel.

One can generate a large voltage using a car battery of 12 volts with the help of inductors. Inductors are known for opposing rapid changes in current, which make them the perfect tool for generating sparks.

In the ignition system of an automobile, a specific type of inductor, referred to as the spark coil, is employed. By creating a significant change in current over a brief period, the voltage across the inductor can be amplified. When the ignition switch is engaged, the current passing through the inductor gradually increases until it reaches a stable state. At this juncture, the rate of current change and the voltage of the inductor are both zero.

However, when the switch is suddenly disengaged, a high voltage is created across the inductor due to the rapidly collapsing magnetic field. This results in a spark or arc in the air gap. The spark persists until all the energy stored in the inductor is exhausted in the spark discharge.

This can be illustrated with an example of designing an automobile ignition system. Assuming that the system's spark coil has an inductance of 20-mH and a resistance, and is supplied with a voltage of 12 V. The task at hand is to calculate the time required for the coil to charge completely, the energy stored in the coil, and the voltage generated at the spark gap if the switch opens in 2 microseconds.

After conducting the necessary calculations, the conclusion drawn is that the coil will be fully charged in 20 milliseconds, will store an energy of 57.6 millijoules, and will generate a voltage of 24 kilovolts at the spark gap when the switch is opened. This high voltage is responsible for creating the spark that ignites the fuel-air mixture in the cylinder, thereby powering the engine of a car.

 Core: Electrical Engineering

An Introduction to the Laboratory Mouse: Mus musculus

JoVE 5129

Mice (Mus musculus) are an important research tool for modeling human disease progression and development in the lab. Despite differences in their size and appearance, mice share a distinct genetic similarity to humans, and their ability to reproduce and mature quickly make them efficient and economical candidate mammals for scientific study.

This video provides a brief overview of mice, both as organisms and in terms of their many advantages as experimental models. The discussion features an introduction to common laboratory mouse strains, including the nude mouse, whose genetic makeup renders them both hairless and immunodeficient. A brief history of mouse research is also offered, ranging from their first use in genetics experiments to Nobel prize-winning discoveries in immunology and neurobiology. Finally, representative examples of the diverse types of research that can be performed in mice are presented, such as classic behavioral tests like the Morris water maze and in-depth investigations of mammalian embryonic development.

 Biology II

Tissue Regeneration with Somatic Stem Cells

JoVE 5339

Somatic or adult stem cells, like embryonic stem cells, are capable of self-renewal but demonstrate a restricted differentiation potential. Nonetheless, these cells are crucial to homeostatic processes and play an important role in tissue repair. By studying and manipulating this cell population, scientist may be able to develop new regenerative therapies for injuries and diseases.

This video first defines somatic stem cells, and then explores the role these cells play in tissue regeneration. This is emphasized in a description of a protocol that isolates muscle satellite cells and uses them to repair muscle damage in a mouse model of muscular dystrophy. Finally, we discuss specific tissue regeneration studies utilizing somatic stem cells.

 Developmental Biology

Anxiety Testing

JoVE 5430

Anxiety is a commonly observed behavioral disorder that stems from fear. It is described as increased restlessness, or unpleasant feelings of fear over anticipated events. Experimenters often use rodent models to better understand anxiety disorders in humans. They use different paradigms, like exposing rodents to bright spaces or loud sounds, which are known to induce fear. These tests combined with other interventions such as surgery or drug-administration may assist researchers in pinpointing the neurobiological basis of anxiety disorders.

This video begins by providing common principles behind variety of anxiety tests. Then, two specific protocols, the Successive Alleys Test and the Hyponeophagia Test are discussed in detail.  Lastly, variations of anxiety testing in rodents and humans will be explored.

 Behavioral Science

Cytogenetics

JoVE 5545

Cytogenetics is the field of study devoted to chromosomes, and involves the direct observation of a cell’s chromosomal number and structure, together known as its karyotype. Many chromosomal abnormalities are associated with disease. Each chromosome in a karyotype can be stained with a variety of dyes to give unique banding patterns. More recent techniques, including comparative genomic hybridization and fluorescence in situ hybridization (FISH), allow for detecting specific chromosomal features or abnormalities.

This video will begin by examining the principles of these classical and modern cytogenetics techniques. This is followed by an examination of a general protocol for performing FISH. Finally, several examples of how karyotyping can be applied to various medical applications are presented.

 Genetics

An Introduction to Cell Division

JoVE 5640

Cell division is the process by which a parent cell divides and gives rise to two or more daughter cells. It is a means of reproduction for single-cell organisms. In multicellular organisms, cell division contributes to growth, development, repair, and the generation of reproductive cells (sperms and eggs). Cell division is a tightly regulated process, and aberrant cell division can cause diseases, notably cancer.

JoVE's Introduction to Cell Division will cover a brief history of the landmark discoveries in the field. We then discuss several key questions and methods, such as cell cycle analysis and live cell imaging. Finally, we showcase some current applications of these techniques in cell division research.

 Cell Biology

Nuclear Magnetic Resonance (NMR) Spectroscopy

JoVE 5680

Source: Laboratory of Dr. Henrik Sundén – Chalmers University of Technology

Nuclear magnetic resonance (NMR) spectroscopy is a vital analysis technique for organic chemists. With the help of NMR, the work in the organic lab has been facilitated tremendously. Not only can it provide information about the structure of a molecule but also determine the content and purity of a sample. Compared with other commonly encountered techniques for organic chemists — such as thermal analysis and mass spectrometry (MS) — NMR is a non-destructive method that is valuable when recovery of the sample is important.

One of the most frequently used NMR techniques for an organic chemist is proton (1H) NMR. The protons present in a molecule will behave differently depending on its surrounding chemical environment, making it possible to elucidate its structure. Moreover, it is possible to monitor the completion of a reaction by comparing NMR spectra of the starting material to that of the final product.

This video exemplifies how NMR spectroscopy can be used in the everyday work of an organic chemist. The following will be shown: i) preparation of an NMR sample. ii) Using 1H NMR to monitor a reaction. iii) Identifying the product obtained from a reaction with 1H NMR. The reaction that will be shown is the synthesis of an E-chalcone (3) from an aldehyde (1) and a ketone (2) (Scheme 1).1

Scheme 1
Scheme 1. Synthesis of (2E)-3-(4-methoxyphenyl)-1-(4-methylphenyl)-2-propen-1-one.

 Organic Chemistry

Introduction to Titration

JoVE 5699

Source: Laboratory of Dr. Yee Nee Tan — Agency for Science, Technology, and Research

Titration is a common technique used to quantitatively determine the unknown concentration of an identified analyte.1-4 It is also called volumetric analysis, as the measurement of volumes is critical in titration. There are many types of titrations based on the types of reactions they exploit. The most common types are acid-base titrations and redox titrations.5-11

In a typical titration process, a standard solution of titrant in a burette is gradually applied to react with an analyte with an unknown concentration in an Erlenmeyer flask. For acid-base titration, a pH indicator is usually added in the analyte solution to indicate the endpoint of titration.12 Instead of adding pH indicators, pH can also be monitored using a pH meter during a titration process and the endpoint is determined graphically from a pH titration curve. The volume of titrant recorded at the endpoint can be used to calculate the concentration of the analyte based on the reaction stoichiometry.

For the acid-base titration presented in this video, the titrant is a standardized sodium hydroxide solution and the analyte is domestic vinegar. Vinegar is an acidic liquid that is frequently used as a culinary condiment or flavorings. Vinegar mainly consists of acetic acid (CH3COOH) and water. The acetic acid content of commercial vinegar can vary widely and the goal of this experiment is to determine the acetic acid content of commercial vinegar by titration.

 General Chemistry

Electrospinning of Silk Biomaterials

JoVE 5798

Silk fibers have been processed and used to create fabrics and threads for centuries. However, the solubilizing of silk fibers, thereby turning it into a versatile pre-polymer solution is a much newer technology. Solubilized silk can be processed in many different ways to create a biocompatible material with controllable mechanical properties.

This video introduces the processing of silk from silk worm cocoons, and shows how the silk solution can be used to create a fiber mat via electrospinning. Several applications of this technique, such as its use as a structural material in engineered tissue scaffolds, are then introduced.

 Bioengineering

Chirality

JoVE 11719

Chirality is a term that describes the lack of mirror symmetry in an object. In other words, chiral objects cannot be superposed on their mirror images. For example, our feet are chiral, as the mirror image of the left foot, the right foot, cannot be superposed on the left foot.

Chiral objects exhibit a sense of handedness when they interact with another chiral object. For example, our left foot can only fit in the left shoe and not in the right shoe. Achiral objects — objects that have superposable mirror images — do not have a sense of handedness. For example, socks are achiral; as such, a sock can be worn on both feet equally well.

Chiral objects are identified by the lack of certain symmetry elements in their structure. Specifically, chiral objects lack a plane of symmetry: an imaginary plane that can divide an object into two equal halves. In addition, chiral objects also lack a center of symmetry: a point from which similar components of the object are equidistant and opposite to each other.

Chiral molecules exist as a pair of non-superposable mirror images. As a rule of thumb, molecules which have only one tetrahedral carbon atom with four different substituents attached to it are always chiral. Such a tetrahedral carbon atom is referred to as a chiral center. In general, to identify whether a molecule is chiral, the molecular geometry should be known. If the molecular geometry lacks a plane of symmetry as well as a center of symmetry, the molecule is chiral.

 Core: Organic Chemistry

Degree of Unsaturation

JoVE 11766

The degree of unsaturation (U), or index of hydrogen deficiency (IHD), is defined as the difference in the number of pairs of hydrogen atoms between the compound and the acyclic alkane with the same number of carbon atoms. Each double bond or ring costs two hydrogen atoms compared to a saturated analog and results in one degree of unsaturation.

The degree of unsaturation for hydrocarbons is U = (2C + 2 − H) / 2, where C is the number of carbon atoms and H is the number of hydrogen atoms.

For substituted hydrocarbons and ions, the degree of unsaturation is calculated by adjusting the equation for the difference in valency.

In organohalogen compounds, the halogen atoms, being monovalent, are treated as equivalent to hydrogen atoms. Thus, the formula to calculate the degree of unsaturation for compounds containing halogen atoms is U = (2C + 2 − (H + X)) / 2, where X is the number of halogen atoms.

Divalent oxygen atoms in organooxygen compounds are ignored while calculating the saturation number. Each trivalent nitrogen atom in organonitrogen compounds adds to the numerator: U = (2C + 2 + N − H) / 2, where N is the number of nitrogen atoms.

Accordingly, the overall formula for degree of unsaturation is U = (2C + 2 + N − (H + X)) / 2.

 Core: Organic Chemistry

Reduction of Alkenes: Catalytic Hydrogenation

JoVE 11785

Alkenes undergo reduction by the addition of molecular hydrogen to give alkanes. Because the process generally occurs in the presence of a transition-metal catalyst, the reaction is called catalytic hydrogenation.

Metals like palladium, platinum, and nickel are commonly used in their solid forms — fine powder on an inert surface. As these catalysts remain insoluble in the reaction mixture, they are referred to as heterogeneous catalysts.

The hydrogenation process takes place on the surface of the metal catalyst. It begins with the adsorption of the hydrogen onto the metal surface, followed by the cleavage of the H–H bonds to give individual metal–hydrogen bonds. The alkene then complexes with the catalyst surface by using its p orbitals to overlap with the empty metal orbitals of the catalyst. The two hydrogen atoms then insert into the π bond sequentially through syn addition (addition to the same face of the π bond) to give the reduced product — the alkane. The alkane formed is no longer bound to the metal and diffuses away from the catalyst's surface.

The process of hydrogenation is exothermic. The heat released is called the heat of hydrogenation (ΔH°), and it helps predict the relative stabilities of alkenes. For example, although the hydrogenation of both cis-2-butene and trans-2-butene gives the same product — butane, trans-2-butene is more stable than cis-2-butene. This can be explained based on the heat of hydrogenation of the two isomers. The cis isomer (ΔH° = −28.6 kcal/mol) has a slightly higher heat of hydrogenation compared to the trans isomer (ΔH° = −27.6 kcal/mol). In cis-2-butene, the steric repulsion between the two methyl groups lying on the same side of the double bond makes it less stable, which is reflected in its larger heat of hydrogenation.

 Core: Organic Chemistry

Multiple Halogenation of Methyl Ketones: Haloform Reaction

JoVE 11843

A method involving the transformation of methyl ketones to carboxylic acids using excess base and halogen is called the haloform reaction. It begins with the deprotonation of α hydrogen to form an enolate ion which reacts with the electrophilic halogen to give an α-halo ketone. The step continues until all the α protons are substituted to form a trihalomethyl ketone. The resulting molecule is unstable, and in the presence of a hydroxide base, it readily undergoes nucleophilic acyl substitution. This leads to the expulsion of trihalomethyl carbanion and produces carboxylic acid. The carbanion generated is stable owing to the electron-withdrawing effect of the three halogens. Subsequent deprotonation of the acid by carbanion forms a carboxylate and haloform, which is the driving force of the reaction. Finally, acidification of the carboxylate gives the desired product, and the reaction is named after the by-product. Using chlorine or bromine results in immiscible liquids of chloroform and bromoform. In contrast, iodine forms a yellow precipitate of iodoform, often used to detect methyl ketones in unknown substrates.

 Core: Organic Chemistry

Alcohols from Carbonyl Compounds: Reduction

JoVE 11925

Reduction is a simple strategy to convert a carbonyl group to a hydroxyl group. The three major pathways to reduce carbonyls to alcohols are catalytic hydrogenation, hydride reduction, and borane reduction.

Catalytic hydrogenation is similar to the reduction of an alkene or alkyne by adding H2 across the pi bond in the presence of transition metal catalysts like Raney Ni, Pd–C, Pt, or Ru. Aldehydes and ketones can be reduced by this method, often under mild to moderate heat (25–100°C) and pressure (1–5 atm H2), to yield 1° and 2° alcohols, respectively.

Figure1

Figure 1. Catalytic hydrogenation can be suitable for industrial applications when harsh conditions are not required, but unsaturated carbon–carbon bonds are also reduced.

Hydride reduction can be achieved by hydride transfer reagents, like sodium borohydride (NaBH4) and lithium aluminum hydride (LiAlH4, or LAH), as nucleophilic attack by a free hydride ion, Figure2, is almost unknown for NaH salts due to its high charge density, making it a strong base. The hydrogen atoms of LAH and NaBH4, being covalently bonded to boron and aluminum atoms, have partial negative charges, thereby enhancing their nucleophilicity at the cost of basicity. The first step of nucleophilic addition leads to the formation of alkoxide ions. The byproduct alkoxyborohydride or alkoxyaluminate reduces three more carbonyl molecules, successively transferring all their hydrogen atoms. Since hydride is a poor leaving group, the hydride transfer steps are irreversible, and therefore the reaction proceeds to completion. Lastly, the reaction mixture is worked up with solvent (i.e., water or alcohol in the case of NaBH4 and dilute acid in the case of LAH).

Figure3

Figure 2. Aldehydes and symmetrical esters produce one 1° alcohol product. Unsymmetrical esters produce a mixture of 1° alcohols.

LAH, NaBH4, and their derivatives are highly useful in the reduction of aldehydes and ketones. LAH, a powerful reducing agent, can also reduce carbonyl compounds like acids, esters, acyl chlorides, and amides. LAH reacts violently with water and other protic solvents, liberating hydrogen gas and forming metal hydroxides/alkoxides. Hence, LAH reductions are typically carried out in aprotic solvents like anhydrous diethyl ether and THF.

Figure4

Figure 3. An alcoholic solution of lithium borohydride is a non-hazardous alternative to LAH in selectively reducing esters over acids.

On the other hand, NaBH4 is of a milder nature and generally reduces only in protic solvents like ethanol or methanol. Therefore, NaBH4 can be used to reduce aldehydes and ketones in the presence of functional groups like alkyl halides, esters, alkyl tosylates, and nitro groups. Diisobutylaluminum hydride (DIBAL-H) can also convert carbonyls to alcohols at room temperature by two successive additions of hydride ions. However, when reacted with esters at low temperatures, this reaction can be stopped at the aldehyde stage by adding only one equivalent of hydride ion.

Borane reduction uses a borane (BH3) solution in diethyl ether, THF, or Me2S to selectively reduce electron-rich carbonyl groups like carboxylic acids in the presence of other reducible functional groups such as esters and even ketones.

Figure5

Figure 4. The formation of a reactive triacylborate intermediate with a more electrophilic carbonyl group than the starting ester molecule drives the reduction reaction forward.

For living organisms, reduced coenzyme NADH or its phosphoester NADPH is equivalent to laboratory hydride reagents in enzyme catalysis of similar biological reductions.

 Core: Organic Chemistry

Formation of Muscle Fibers from Myoblasts

JoVE 12522

De novo myogenesis, or the formation of muscle fibers, begins during the early embryonic stages. The skeletal muscle is formed from somites– blocks of embryonic cell layers. The somites are further divided into dermatomes, myotomes, sclerotomes, and syndetomes. Among these, the myotomes give rise to muscle fibers.

Muscle progenitor cells (MPCs) are formed from the myotomes. MPCs express genes that encode the transcription factors Pax3 and Pax7. Along with Pax 3/7, other transcription factors from the myogenic regulatory factors drive the MPCs into the myogenic lineage. Pax3 activates the expression of the MyoD–a basic helix-loop-helix transcription factor and a strong transactivator. MyoD can transform MPCs and undifferentiated cells, including fibroblasts, into muscle cells such as myoblasts through myogenesis. The dividing myoblasts exit the cell cycle and express genes that transform myoblasts into myofibrils.

Mitochondria play an important role in myogenesis. During myogenesis, cells switch from primarily glycolysis for energy to oxidative phosphorylation to produce the ATP required for myogenesis. Oxidative phosphorylation takes place in the mitochondrial matrix. Additionally, mitochondrial enzymes are highly active during myoblast differentiation and muscle regeneration. Similarly, muscle injury causes a loss of citrate synthase activity, resulting in aberrant muscle regeneration. As the regulatory role of mitochondria in myogenesis is being investigated, mitochondria appear to be promising candidates for the treatment of muscle-related disorders.

 Core: Cell Biology

Tonicity in Plants

JoVE 13366

Plant cells maintain appropriate osmotic balance in extreme conditions. For instance, plants in dry environments store water in vacuoles, limit the opening of their stoma, and have thick, waxy cuticles to prevent unnecessary water loss. Some species of plants that live in salty environments store salt in their roots. As a result, water osmosis occurs in the root from the surrounding soil.

Tonicity

Tonicity describes the capacity of a cell to lose or gain water depending on the solute concentration outside. Organisms such as plants, fungi, bacteria, and protists, have cell walls surrounding the plasma membrane. Three possible scenarios alter the volume of a cell: hypertonicity, hypotonicity, and isotonicity.

Hypotonic environment

In hypotonic environments, there is a higher concentration of solutes inside plant cells than outside. Water enters the cell via osmosis and causes it to swell. Because the cell wall limits the expanding plasma membrane, the cell does not lyse. By limiting expansion, the cell wall allows cells to become turgid, resulting in the stiffening of plants.

Plant cytoplasm is always slightly hypertonic to the cellular environment, and water will always enter a cell if water is available. The force generated when an influx of water causes the plasma membrane to push against the cell wall is called turgor pressure. Turgor pressure keeps non-woody plants upright.

Hypertonic Environments

Conversely, the extracellular fluid becomes hypertonic in a dry climate, causing water to leave the cell through osmosis. In this condition, the cell cannot shrink because the cell wall is not flexible. As a result, vacuoles decrease in size the cell membrane detaches from the wall and constricts the cytoplasm. This process is called plasmolysis. Thus, plants lose turgor pressure and wilt.

 Core: Cell Biology

Enzyme-Linked Immunosorbent Assay

JoVE 13383

In 1971, Peter Perlman and Eva Engvall developed an Enzyme-linked immunosorbent assay (ELISA or EIA). ELISA differs from western blot in that the assays are conducted in microtiter plates or in vivo rather than on an absorbent membrane.

There are many different types of ELISAs, but they all involve an antibody molecule whose constant region binds an enzyme, leaving the variable region free to bind its specific antigen.  Enzyme-substrate reaction allows the antigen to be visualized or quantified. In ELISAs, the substrate for the enzyme is most often a chromogen, a colorless molecule that is converted into a colored end product. The most widely used enzymes are alkaline phosphatase and horseradish peroxidase, for which appropriate substrates are readily available. In some ELISAs, the substrate is a fluorogen, a nonfluorescent molecule that the enzyme converts into a fluorescent form. ELISAs that utilize a fluorogen are  termed  fluorescent enzyme immunoassays (FEIAs). Fluorescence can be detected by either a fluorescence microscope or a spectrophotometer.

There are several types of ELISA, based on differences in the format of detection and general workflow. In direct ELISA, antigens are immobilized in the well of a microtiter plate. An antibody that is specific for a particular antigen, and is conjugated to an enzyme, is added to each well. If the antigen is present, then the antibody will bind. After washing to remove any unbound antibodies, a chromogen is added.The presence of the enzyme converts the substrate into a colored end product.

Indirect ELISA is an extremely sensitive and flexible procedure. Here, in addition to a primary antibody, a secondary antibody is added for detection purposes. The secondary antibody  quantifies how much antigen-specific antibody is present in the sample by the intensity of the color produced from the conjugated enzyme-chromogen reaction.

Sandwich ELISA is more specific and sensitive than direct and indirect ELISA. The goal is to use antibodies to precisely quantify the specific antigen present in a solution, such as the antigen from a pathogen, a serum protein, or hormone from blood or urine. The primary antibody captures the antigen and, following a wash, the polyclonal enzyme-conjugated secondary antibody is added. After a final wash, a chromogen is added, and the enzyme converts it into a colored end product. The amount of color produced (measured as absorbance) is directly proportional to the amount of enzyme, which in turn is directly proportional to the captured antigen. The complex workflow and several optimizations make this process error-prone.

For competitive ELISA, crude samples can be directly used. The experimental setup is highly flexible, wherein direct, indirect, or sandwich ELISAs can be adapted to a competitive format. The sample analyte concentration is determined by the signal interference. The "competition" comes from the fact that if a sample antigen is being tested, then the incubation of the sample with a primary antibody will result in lesser antibodies available to bind to the wells coated with the same antigen. Thus, the intensity of the signal produced in the well, due to the competitive binding, which is concentration dependent, becomes inversely correlated to the amount of sample antigen.

Depending on the purpose, the subtype of ELISA is chosen for a suitable outcome. For detecting large proteins comprising multiple epitopes, sandwich ELISA is most appropriate, whereas competitive ELISA is ideal for small protein  detection. ELISA has significant applications in diagnostics, such as pregnancy testing, detecting food allergens, and identifying cancer biomarkers.

 Core: Cell Biology

Overview of Electron Microscopy

JoVE 13399

The wavelengths of visible light ultimately limit the maximum theoretical resolution of images created by light microscopes. Most light microscopes can only magnify 1000X, and a few can magnify up to 1500X. Electrons, like electromagnetic radiation, can behave like waves, but with wavelengths of 0.005 nm, they produce significantly greater resolution up to 0.05 nm as compared to 500 nm for visible light. An electron microscope (EM) can create a sharp image that is magnified up to 2,000,000X. Thus, EMs can resolve subcellular structures and some molecular structures (for example, single strands of DNA).

There are two basic types of EM — the transmission electron microscope (TEM) and the scanning electron microscope (SEM). The TEM is somewhat analogous to the brightfield light microscope in terms of the way it functions. However, it uses an electron beam focused on the sample from above using a magnetic lens (rather than a glass lens) and projected through the sample onto a detector. Electrons pass through the specimen, and then the detector captures the image.

For electrons to pass through the specimen in a TEM, the sample must be extremely thin (20–100 nm thick). The image is produced because of varying opacity in various parts of the sample. This opacity can be enhanced by staining the specimen with materials such as heavy metals, which are electron-dense. TEM requires that the beam and sample be in a vacuum and that the sample be ultrathin and dehydrated.

SEMs form images of surfaces of specimens, usually from electrons that are knocked off the sample by a beam of electrons. This can create highly detailed images with a three-dimensional appearance displayed on a monitor. Typically, specimens are dried and prepared with fixatives that reduce artifacts before being sputter-coated with a thin layer of metal such as gold. Whereas transmission electron microscopy requires very thin sections and allows one to see internal structures such as organelles and the interior of membranes, scanning electron microscopy can be used to view the surfaces of larger objects (such as a pollen grain) as well as the surfaces of very small samples.

This text is adapted from Openstax, Microbiology, Section 2.3: Instruments of Microscopy

 Core: Cell Biology

Stem Cell Niche

JoVE 13463

The stem cell niche is the dynamic microenvironment where stem cells reside. Inside these niches, the cells may remain undifferentiated, undergo high self-renewal, or become lineage-specific progenitors. Stem cells coexist with other niche cells, such as stromal cells. They also interact closely with the ECM. Cell-cell and cell-matrix communication occur via adhesion molecules or soluble factors that signal the stem cells and determine their fate. Stromal cells also provide survival signals to the stem cells preventing their apoptotic death. This way, the niche allows stem cells to produce progenitors or transit-amplifying cells (TA cells) periodically and replace the body’s damaged or dead cells. Thus, the niche maintains a balance between stem cell quiescence and differentiation.

Adult tissues, including the bone marrow, skin, intestine, or brain, harbor stem cells inside specific niches. For example, hematopoietic stem cells (HSCs) reside amongst osteoblastic cells, stromal cells, and reticulocytes that form the bone marrow niche. The epithelial stem cells of the skin live in the bulge area of the hair follicles. These stem cells interact closely with keratinocytes and help regenerate the hair follicles. The neural stem cells of the adult nervous system are found within the hippocampus region that produces neuroblasts and mature neurons. The neighboring endothelial cells in the hippocampus form the stem cell niche of the nervous system. In the intestine, the intestinal stem cells or ISCs are found in the crypt region interspersed with the Paneth cells. Paneth cells constitute the intestinal stem cell niche and induce the ISCs to produce TA cells and replace the villus every 3 to 5 days.

 Core: Cell Biology

Maintenance of the ES Cell State

JoVE 13480

The cells of the blastocyst inner cell mass only remain pluripotent for a short time. This state of pluripotency and self-renewal can be maintained in embryonic stem (ES) cell culture by adding specific chemicals or growth factors to ensure the cells can continue dividing and later differentiate into different cell types. In some cases, the cells are grown on a feeder layer of differentiated cells, which provides the growth factors and extracellular matrix components necessary for stem cell proliferation.

The ES cell state is regulated by the transcription factors, Oct4, Sox2, and Nanog. These factors activate the genes required for maintaining pluripotency and repress those involved in differentiation. Silencing any of these transcription factors results in lineage-specific differentiation. The factors operate along with many other transcription factors, including Klf4, Klf5, and Smad1, and transcriptional cofactors, such as p300, Mediator, and Nipb, help activate or repress genes without directly binding to DNA.

Chromatin regulators play a crucial role in maintaining the ES cell state. These are categorized as histone-modifying enzymes and ATP-dependent chromatin regulators. Histone modifying enzymes alter the DNA-histone interaction and transcriptionally activate specific genes. These enzymes also suppress development regulators such as polycomb group (PCG) protein complexes. ATP-dependent chromatin remodeling complexes use the energy released by ATP hydrolysis to either displace histones from DNA or enable the relocation of histones. This weakens the bonds between histones and DNA, providing access to transcription factors to bind to DNA.

 Core: Cell Biology

Tissue Homogenization and Cell Lysis

JoVE 13830

Tissue homogenization involves disintegrating tissue architecture and lysing cells, and is an early step in isolating and analyzing cellular components. The method used for homogenization depends on the sample type, the amount of sample available, the analyte to be obtained, and the sensitivity of the method. These methods are broadly classified as mechanical and non-mechanical methods.

Mechanical methods of tissue homogenization

These methods rely on applying external physical force to disrupt tissues and cells. They make use of specialized tools and instruments for homogenization. These instruments use grinding, shearing, blending, beating, or shock to disintegrate the sample. For example, in a French press, the sample is pushed through a small opening under pressure which causes the cells to disrupt. Other homogenizers such as Waring blenders and rotor-stators cut and shear the tissues into significantly smaller sizes.

Non-mechanical methods of tissue homogenization

Non-mechanical lysis methods involve chemical disruptions rather than physical forces to lyse the cells. The tissue is homogenized in a lysis buffer that regulates pH, ionic strength, osmotic strength, and enzymatic activity. The lysis buffer thus aids cell lysis and protects the cell components from damage.

While the enzymes of the lysis buffer help degrade the extracellular matrix of tissues to release individual cells, the surfactant or detergent helps disrupt cell membranes and denature proteins. Sodium dodecyl sulfate (SDS) and Triton-X 100 are two popularly used detergents in these buffers. Another component, the chaotropes, disrupt weak interactions between molecules, thus denaturing the proteins and keeping nucleic acids intact during isolation.

Other non-mechanical physical methods involve using temperature cycles in which the sample is frozen on dry ice or in an ethanol bath and then thawed at room temperature or 37℃. These repeated cycles cause the cell membranes to weaken and rupture. Cell membranes can also be ruptured by osmotic imbalance by placing them in a hypotonic or hypertonic solution. The inward or outward movement of water due to the osmotic gradient causes the cells to swell and burst, or shrink and collapse, releasing their internal contents.

While numerous methods and tools are available for homogenization, each has pros and cons that must be evaluated based on the specific requirements. Often, mechanical methods alone may not wholly or efficiently homogenize a sample. In such cases, mechanical methods may be combined with non-mechanical methods for complete homogenization.

While numerous methods and tools are available for homogenization, each has pros and cons that must be evaluated based on the specific requirements. Often, mechanical methods alone may not wholly or efficiently homogenize a sample. In such cases, mechanical methods may be combined with non-mechanical methods for complete homogenization.

 Core: Cell Biology

Uncertainty: Confidence Intervals

JoVE 14513

The confidence interval is the range of values around the mean that contains the true mean. It is expressed as a probability percentage. The interpretation of a 95% confidence interval, for instance, is that the statistician is 95% confident that the true mean falls within the interval. The upper and lower limits of this range are known as confidence limits. The confidence limits for the true mean are estimated from the sample's mean, the standard deviation, and the statistical factor 't,' or t-score, which depends on the number of degrees of freedom and the desired confidence level. It is important to specify whether a one- or two-tailed confidence interval is needed because the confidence level and the one-tailed t-score table differs from the two-tailed version. As the number of measurements increases, the deviation from the mean decreases, leading to a narrow confidence interval.

 Core: Analytical Chemistry

Chemical Equilibria: Systematic Approach to Equilibrium Calculations

JoVE 14529

Equilibrium calculations for systems involving multiple equilibria are often complex. For example, to calculate the solubility of a sparingly soluble salt in an aqueous solution in the presence of a common ion, one must consider all the equilibria in this solution. Calculations for these systems can be complicated and tedious, so a systematic approach with a series of steps is often helpful. The process is detailed below.

The first step is to identify all the chemical reactions involved, The next is to formulate their corresponding equilibrium constant expressions. The third step is to write the equilibrium mass balance equations, which are based on the law of conservation of mass. These equations relate the total amount of a species in the solution in different forms (equilibrium concentration of the species) to the total amount of that species initially added to the solution (analytical concentration of the species). There may be more than one mass-balance equation for a given set of equilibria. Furthermore, for equilibria with charged species, a charge balance equation is required based on the principle of electroneutrality, which states that the total charge concentration of the positively charged species in the solution must equal that of the negatively charged species. For a given set of equilibria, only one charge-balance equation can be written. After mass-balance and charge-balance equations are written, the number of independent equations obtained is counted. If it equals or exceeds the number of unknowns involved in the equilibria, the equations can be solved for the unknowns. After solving for the unknowns, the validity of the approximations needs to be checked because approximations of the relative concentrations of species are often introduced in the mass-balance equations or charge-balance equations to simplify the calculations.   

 Core: Analytical Chemistry

Titration of Polyprotic Base with a Strong Acid

JoVE 14545

The titration of a polyprotic base such as sodium carbonate with a strong acid such as hydrochloric acid results in two equivalence points on the titration curve. At the first equivalence point, the carbonate ions in the base are completely converted to bicarbonate ions. The second equivalence point corresponds to the complete conversion of bicarbonate ions to carbonic acid, which dissociates into carbon dioxide and water. The region before the first equivalence point corresponds to the carbonate/bicarbonate buffer, and the area between the first and second equivalence points corresponds to the bicarbonate/carbonic acid buffer. While phenolphthalein is used to detect the first endpoint, a mixture of methyl orange and xylene cyanol is used to sharpen the second endpoint.

 Core: Analytical Chemistry

Precipitate Formation and Particle Size Control

JoVE 14584

In precipitation gravimetry, the precipitating agent should react specifically or selectively with the analyte. While a specific reagent reacts with the analyte alone, a selective reagent can react with a limited number of chemical species.

The obtained precipitate should be either a pure substance of known composition or easily converted to one by a simple process, such as ignition or drying. In addition, the precipitate should be insoluble and easily filterable. In general, filterability increases with the size of the precipitate particles. Colloidal suspensions contain minuscule particles with diameters varying from 10−9 to 10−6 m, which are invisible to the naked eye and not easily filtered. However, crystalline suspensions have larger particles that settle quickly and are readily filtered.

The temperature, precipitate solubility, reactant concentrations, and speed of mixing of reactants can affect the particle size. The overall effect of these attributes is called relative supersaturation, RSS, which can be expressed in terms of the concentration of the solute (Q) and its equilibrium solubility (S). The size of the obtained particles is inversely proportional to the average relative supersaturation when the reagent is added. As a result, when the relative supersaturation ratio is high, colloidal precipitates are favored, while crystalline precipitates with large particle sizes are obtained at low relative supersaturation ratios. 

 Core: Analytical Chemistry

Sublimation

JoVE 14617

Sublimation is the direct transformation of a solid to a gaseous state. For instance, at standard pressure and room temperature, solid carbon dioxide sublimes to gaseous carbon dioxide. The phase diagram depicts the conditions required for sublimation. This process occurs at the solid-gas phase boundary and is not observed above the triple point of the substance. The reverse of sublimation is called deposition, where a gaseous substance condenses directly into a solid. Sublimation and deposition can both be used to separate an analyte from interferents.

In simple sublimation, a solid sample is heated in a beaker covered with an inverted watch glass acting as the secondary surface. The analyte sublimes into a gas, and then cools to a solid again when it meets the watch glass. For successful purification, the sublimation temperature must be high enough to ensure high vapor pressure, but low enough to avoid either melting or decomposing the substance. Iodine is an example of a substance amenable to purification by simple sublimation.

If the sample is sensitive to heat, low-temperature sublimation techniques such as freeze-drying (also known as lyophilization) can be used. These techniques are often applied to dehydrate the desired material. Here, the sample is frozen using an ultra-low-temperature mixture (dry ice/acetone), and placed in a vessel attached to a vacuum pump. The ice in the sample sublimes to water, leaving behind the desired, cryo-desiccated product. Lyophilization is used extensively in purifying enzymes for use in biochemistry, molecular biology, and food preservation. 

 Core: Analytical Chemistry

Muscle Recovery and Fatigue

JoVE 14845

Muscle fatigue refers to the decline in a muscle's ability to maintain the force of contraction after prolonged activity. It primarily stems from changes within muscle fibers. Even before experiencing muscle fatigue, one may feel tired and have the urge to stop the activity. This response, known as central fatigue, occurs due to changes in the central nervous system, namely the brain and spinal cord. While there is no single mechanism that induces fatigue, it may serve as a protective response to prevent muscle damage that can be caused due to overuse.

During intense or prolonged physical activity, the ATP demand in the muscles increases. As contractions become more vigorous, compressed blood vessels hinder the oxygen supply to the tissues. Consequently, muscles resort to anaerobic glycolysis to produce ATP, leading to the production of lactic acid. The accumulation of lactate and hydrogen ions from ATP hydrolysis adversely affects cellular pH, resulting in lactic acidosis, which hampers muscle contraction and leads to fatigue. To recover, muscles require time and ample oxygen to regenerate ATP. Once oxygen levels are restored, lactate dehydrogenase enzymatically converts accumulated lactate back into pyruvate. Mitochondria then utilize pyruvate to produce more ATP or enzymatically convert it into glycogen for storage. The surplus ATP is utilized by the enzyme creatine kinase to replenish phosphocreatine reserves.

In addition, the muscles redirect any excess lactate to the liver. Once in the liver, the lactate converts back into glucose, which is then transported back to the muscles. This cyclic process, known as the Cori cycle, is a vital metabolic pathway that plays a crucial role in the transport and conversion of lactate. It ensures the efficient metabolism of lactate, thereby enabling sustained energy production in the body.

 Core: Anatomy and Physiology

Muscles for Facial Expressions

JoVE 14866

The craniofacial muscles are a collection of approximately 20 thin skeletal muscles situated beneath the skin of the face and scalp. These muscles, primarily responsible for the vast array of human facial expressions, originate from the bones or fibrous structures of the skull and extend outwards to connect with the skin. While most skeletal muscles in the body are enveloped in thick fascia, facial muscles generally have a more delicate fascial covering, with the buccinator muscle being a notable exception.

These facial muscles are situated around the facial openings (mouth, eyes, nose, and ears) or span across the skull and neck and are crucial for non-verbal communication. They are aptly named muscles of facial expression or mimetic muscles due to their specific role in reflecting emotions and intentions. All facial muscles receive signals from the facial nerve, ensuring precise and coordinated movements, while the facial artery provides the necessary blood supply. However, the innervation and vascularization patterns can be complex and vary slightly between individuals.

For instance, the occipitofrontalis muscle, which plays a key role in lifting the eyebrows, is a two-part scalp muscle composed of frontal and occipital bellies. These parts are united by the epicranial aponeurosis, a tough layer of connective tissue. Another crucial muscle is the orbicularis oris, a circular muscle that encircles the mouth. It originates from adjacent muscles attached to the maxillae and mandible and inserts around the lips. It is also the insertion point for the buccinator muscle, which originates from the maxillae and mandible and controls mouth movements like pouting and puckering.

Additionally, muscles such as the zygomaticus major, responsible for pulling the corners of the mouth upward into a smile, and the corrugator supercilii, involved in frowning by drawing the eyebrows together, further exemplify the diversity and specificity of facial muscles. Along with others controlling the eyes, nose, and cheeks, these muscles orchestrate a symphony of expressions ranging from joy to sorrow.

 Core: Anatomy and Physiology

Functional Divisions of the Nervous System

JoVE 14883

The nervous system, responsible for sensing, integrating, and responding to various stimuli, is divided into the central nervous system (CNS) and the peripheral nervous system (PNS). The PNS has two functional divisions: the sensory or afferent division and the motor or efferent division.

The sensory division transmits information from sensory receptors in the body to the CNS. It provides the CNS with knowledge about somatic senses (such as tactile, thermal, pain, and proprioceptive sensations) and special senses (like smell, taste, vision, hearing, and equilibrium).

On the other hand, the motor division carries output from the CNS to effectors like muscles and glands. This division further splits into the somatic nervous system (SNS) and the autonomic nervous system (ANS).

The SNS directs output from the CNS to skeletal muscles. The ANS controls output from the CNS to smooth muscles, cardiac muscles, and glands. The ANS comprises two main branches: the sympathetic nervous system and the parasympathetic nervous system. These usually have opposing actions. The parasympathetic nervous system manages "rest-and-digest" activities, while the sympathetic nervous system supports exercise or emergency actions, also known as "fight-or-flight" responses.

Lastly, there is the enteric nervous system (ENS), a network of over 100 million neurons confined to the wall of the gastrointestinal tract. The ENS helps regulate the activity of the smooth muscle and glands of the GI tract and communicates with and is regulated by the other branches of the ANS.

 Core: Anatomy and Physiology

Neural Circuits

JoVE 14900

Neural circuits and neuronal pools are two of the main structures found in the nervous system. Neural circuits are networks of neurons that work together to carry out a specific task or process. They consist of interconnected neurons and glial cells, which provide structural and metabolic support.

Neuronal pools are collections of nerve cells with similar functions and interact through chemical and electrical signals. These pools include both interneurons (the central neural circuit nodes that enable communication between sensory or motor neurons and the CNS) and projection neurons, which relay information from one part of the brain to another. Neuronal pools can be further classified into ascending pathways (which move information from lower brain regions to higher cortical regions) and descending pathways (which move information from higher cortical regions to lower brain regions).

An example of a neural circuit is the motor cortex, which controls muscle movement. This neural circuit consists of neurons connecting the spine to different parts of the brain and other neuronal pools in the hypothalamus and cerebellum. Similarly, an example of a neuronal pool is the thalamus, which acts as a relay station between different parts of the brain.

Neural circuits and neuronal pools are essential for communication between neurons within the nervous system and for sending signals throughout the body. The combination of these structures can produce complex behaviors in animals or humans by allowing them to remember past experiences and react appropriately in any given situation. In addition, these structures also contribute to learning, as new neural pathways are formed when a person encounters new experiences and information. Ultimately, neural circuits and neuronal pools work together to allow the brain to function normally and carry out daily tasks.

 Core: Anatomy and Physiology

Functional Brain Systems: Reticular Formation

JoVE 14916

The reticular formation is a complex network of gray and white matter located within the brainstem extending from the medulla to the midbrain.

Within the reticular formation, there are several distinct nuclei that can be classified into three broad categories. The Raphe nuclei are located along the midline of the brainstem. They are primarily known for their role in synthesizing and releasing serotonin, a neurotransmitter involved in regulating mood, appetite, sleep, and circadian rhythms. The medial nuclei are located toward the center or middle portion of the reticular formation. They play roles in regulating autonomic functions, such as cardiovascular and respiratory control, as well as in modulating consciousness and arousal. The lateral nuclei are situated toward the outer edges of the reticular formation. They are involved in sensorimotor integration, including regulating muscle tone, reflex activity, attention, and sensory information processing.

Reticular Activating System

The reticular activating system (RAS) lies within the core of the reticular formation. This network of neurons has a complex structure that extends through the core of the brainstem. It spans from the medulla oblongata to the midbrain and regulates various physiological functions, including sleep-wake cycles, arousal, attention, and consciousness. By projecting to the cerebral cortex and through the thalamus, the RAS ensures that the brain remains alert and responsive to external stimuli. Its activation is essential for maintaining a state of wakefulness, while its inactivation promotes sleep onset. Damage to the RAS can result in severe alterations in consciousness, such as coma.

The RAS also heightens attention and alertness, enabling individuals to focus on specific tasks—such as the careful attention required when chopping a cucumber—by processing visual and auditory information. However, it cannot process olfactory stimuli, highlighting the importance of auditory and visual alarms in alerting sleeping individuals to dangers like smoke or fire.

Motor Component

While the RAS performs sensory processing, the motor component of the reticular formation is instrumental in coordinating coarse limb movements and regulating visceral motor functions. Its regulation of muscle tone and reflexes enables the body to adapt to changes in environment and activity, ensuring stability and fluid motion. Additionally, its involvement in visceral motor regulation underpins the autonomic functions critical for life, such as heart rate, breathing patterns, and digestion. By orchestrating these involuntary processes, the reticular formation ensures the body's internal environment is kept in balance, allowing for optimal performance and response to external stresses.

 Core: Anatomy and Physiology

Overview of Somatic Sensory Pathways

JoVE 14934

Somatic sensory or somatosensory pathways refer to the neural pathways that carry information related to touch, pressure, pain, temperature, and proprioception from the skin, muscles, tendons, and joints to the brain. These pathways involve several stages of processing and integration of sensory information.

The somatosensory system is divided into three main pathways: the dorsal (or posterior) column-medial lemniscus, spinothalamic (or anterolateral), and spinocerebellar pathways.

The dorsal column-medial lemniscus pathway transmits information related to proprioception, fine touch, and vibration sensation. It starts with the primary sensory neurons in the spinal cord, which synapse with secondary sensory neurons in the dorsal column nuclei in the medulla oblongata. The axons of the secondary neurons then cross over to the opposite side of the brainstem and ascend through the medial lemniscus to the thalamus, where they synapse with tertiary sensory neurons. The tertiary sensory neurons then project to the primary somatosensory cortex, where the information is processed and integrated.

The spinothalamic pathway transmits information related to pain, temperature, and crude and coarse touch sensations. It starts with the primary sensory neurons, which conduct information from the periphery to the spinal cord. These neurons then synapse with secondary sensory neurons, which cross over to the opposite side of the spinal cord and ascend through the spinothalamic pathway to the thalamus, where they synapse with tertiary sensory neurons. The tertiary sensory neurons then project to the primary somatosensory cortex, where the information is processed and integrated.

The spinocerebellar pathway is a neural pathway that carries sensory information related to proprioception from the muscles and joints to the cerebellum that overall helps in motor coordination and control. Unlike other somatosensory pathways, the spinocerebellar pathway does not involve conscious perception of sensory information.

 Core: Anatomy and Physiology

Sympathetic Signaling

JoVE 14951

Sympathetic signaling, a vital part of the autonomic nervous system, plays a crucial role in mobilizing the body's resources in response to stress or emergencies. It involves the transmission of nerve impulses from sympathetic preganglionic fibers to postganglionic fibers. This results in the release of specific neurotransmitters and activation of adrenergic receptors.

Sympathetic preganglionic fibers release the neurotransmitter acetylcholine (ACh) onto the ganglionic neurons in the sympathetic ganglia. This ACh binds to nicotinic cholinergic receptors on the postganglionic neurons, transmitting nerve impulses. Once activated, the postganglionic fibers release norepinephrine (NE) as the primary neurotransmitter onto their target tissues.

Adrenergic Receptors

The target tissues contain adrenergic receptors, which are classified into two main types: alpha and beta receptors. These receptors are further subdivided into specific subtypes, each with different functional effects upon activation.

Alpha-Adrenergic Receptors

Alpha-adrenergic receptors are divided into alpha-1 and alpha-2 subtypes. Alpha-1 receptors are typically found in smooth muscle cells of blood vessels, causing vasoconstriction when activated. Alpha-2 receptors are present on preganglionic neurons and exert an inhibitory effect. They are also found on the presynaptic terminal of postganglionic neurons, where they function as autoreceptors, regulating NE release. The alpha-2 receptors help coordinate sympathetic and parasympathetic activities.

Beta-Adrenergic Receptors

Beta-adrenergic receptors consist of beta-1, beta-2, and beta-3 subtypes. Beta-1 receptors are primarily found in the heart, where their activation increases heart rate and the force of contraction. Beta-2 receptors are present in smooth muscles lining the airways and cause bronchodilation. Beta-3 receptors are primarily located in adipose tissue, where their activation increases lipolysis and the breakdown of fat cells.

 Core: Anatomy and Physiology

Auditory Pathway

JoVE 14970

Auditory pathways constitute the complex neural circuits responsible for transmitting and interpreting auditory information from the peripheral auditory system to the brain. Sound waves are initially captured by the outer ear, funneled through the ear canal, and reach the tympanic membrane (eardrum). These vibrations are transmitted via the middle ear's ossicles to the inner ear's cochlea.

When viewed cross-sectionally, the cochlea reveals the scala vestibuli and scala tympani flanking the cochlear duct. Numerous Corti organs reside within the cochlear duct, which converts the scala's wave motion into neural impulses. Positioned atop the basilar membrane - which resides between the Corti organs and the scala tympani within the cochlear duct - these organs respond to fluid waves traveling through the scala vestibuli and scala tympani. Locations on the basilar membrane react selectively to wave frequencies; areas proximal to the cochlea base respond to higher frequencies, and areas closer to the cochlea tip react to lower frequencies.

Interspersed within the Corti organs are hair cells, christened for the stereocilia (resembling hair) that project from their apical surfaces. These stereocilia, organized in a gradient from tallest to shortest, are interconnected by protein fibers within each array. These protein tethers facilitate the collective bending of these arrays in response to basilar membrane movement. Extending towards the tectorial membrane - which is affixed medially to the Corti organ - these stereocilia undergo lateral movement as pressure waves from the scala stimulate the basilar membrane. The bending of stereocilia either towards or away from the tallest in the array causes a shift in protein tether tension, opening ion channels within the hair cell membrane if bent towards the tallest and closing them if bent towards the shortest. In the absence of sound, standing stereocilia exert a small degree of tension on the tethers, resulting in a slight depolarization of the hair cell membrane.

The hair cells convert mechanical vibrations into electrical signals, activating the auditory nerve fibers. These signals travel through the auditory nerve to the brainstem, specifically the cochlear nuclei, and ascend through multiple relays, including the superior olivary complex and the inferior colliculus.

The auditory signals continue their journey to the thalamus and ultimately arrive at the auditory cortex in the brain's temporal lobe. This region processes the information, distinguishing various sound attributes such as pitch, intensity, and localization, enabling the perception and interpretation of auditory stimuli.

 Core: Anatomy and Physiology

Anatomy of the Adrenal Glands

JoVE 14986

The adrenal or supra-renal glands, situated above the kidneys and aligned with the twelfth rib, are paired pyramid-shaped structures crucial for the body's stress response. During stress, these glands secrete hormones vital for adaptive physiological reactions.

These glands possess a distinctive yellow tinge due to the stored cholesterol and fatty acids required for hormone synthesis. They are encased in a fibrous capsule and cushioned by fat.

The adrenal gland comprises two distinct regions with specialized functions. The outer layer, the adrenal cortex, constitutes most of the gland and synthesizes corticosteroids. These hormones produced by the zona glomerulosa, zona fasciculata, and zona reticularis regions facilitate prolonged responses to stress.

Encapsulated by the cortex, the interior layer is the adrenal medulla, composed of nervous tissue. This region houses spherical chromaffin cells surrounding capillaries. Chromaffin cells play a pivotal role in synthesizing catecholamines, such as adrenaline and noradrenaline. These hormones induce immediate and brief effects, precipitating the well-known fight-or-flight response during stressful situations.

In summary, the adrenal glands are crucial in orchestrating both immediate and protracted physiological responses to stress, ensuring the body's adaptability and survival in challenging situations.

 Core: Anatomy and Physiology

Nodal Analysis with Voltage Sources

JoVE 15048

Nodal analysis is a remarkably effective method used in electrical engineering to simplify the analysis of complex circuits, including those with dependent or independent voltage sources. Its strength lies in its systematic approach to breaking down circuits into manageable components, making it easier for engineers to understand and solve.

Consider a circuit that contains four resistors and two voltage sources, as shown in Figure 1. One of these voltage sources is connected between a non-reference node and the reference node. In this configuration, the voltage at the non-reference node can be written directly as equal to the voltage of the source. This simplifies the problem by reducing the number of unknowns.

Figure1

The other voltage source is connected between two non-reference nodes. This configuration forms what is known as a supernode or generalized node. Supernodes are a unique concept in nodal analysis, which helps to handle situations where voltage sources are connected between non-reference nodes.

Analyzing a circuit with a supernode necessitates the application of both Kirchhoff's Current Law (KCL) and Kirchhoff's Voltage Law (KVL). First, KCL is applied to the supernode. This law states that the sum of currents entering a node equals the sum of currents leaving it. By considering the currents through each element within the supernode, an equation can be obtained and written in terms of the node voltages.

Next, to apply KVL, the circuit is redrawn to highlight the loop containing the supernode. KVL states that the sum of the electromotive forces in any closed loop or mesh in a network is equal to the sum of the potential drops in that loop. Going around this loop in a clockwise direction gives a constraint equation.

With these steps, three equations are obtained, which represent the behavior of the circuit. These equations can be simultaneously solved to determine the node voltages.

 Core: Electrical Engineering

Kirchhoff's Voltage Law

JoVE 15071

Kirchhoff's Voltage Law (KVL) is another fundamental principle in electrical engineering, introduced by physicist Gustav Robert Kirchhoff. This law is rooted in the principle of energy conservation, which states that energy can neither be created nor destroyed, only transferred or converted from one form to another.

KVL states that the algebraic sum of all voltages around a closed path or loop within a circuit is zero. This means that the total voltage supplied in a loop is equal to the total voltage drop across the components in that loop.

The direction of the loop around the circuit could be clockwise or anticlockwise and can start from any point. Once the loop direction is chosen, positive voltages are those where the negative terminal is encountered first, and negative voltages are those where the positive terminal is encountered first.

By applying KVL in the clockwise direction and rearranging the terms, there is an alternative form of the law: the sum of voltage drops across the components in a loop equals the sum of supplied voltages.

Consider, for example, a household lighting system where three LED lights connected in series are strung together. Each light requires three volts to turn on. By applying KVL, we can determine the voltage needed for the battery to power these lights. The battery voltage equals the sum of the voltage drop across the three lights, which is nine volts.

However, for KVL to hold true, a crucial condition must be met: a circuit cannot contain two different voltages in parallel unless they are equal. This is because, in a parallel connection, the same voltage is applied across all components, and any discrepancy would violate the law of energy conservation.

 Core: Electrical Engineering

Equivalent Capacitance

JoVE 15087

From the study of resistive circuits, it is understood that employing a series-parallel combination serves as an effective strategy for simplifying circuits. Capacitors can be arranged within a circuit in one of two ways: a series configuration or a parallel configuration. The way these capacitors are connected to a battery will influence both the potential drop across each individual capacitor and the size of the charge that each capacitor can store. This is determined by the specific type of connection in place. To simplify this scenario, the combination of capacitors can be substituted with a single equivalent capacitor. This equivalent capacitor is able to store an identical amount of charge as the original combination when subjected to the same potential difference.

When there is a parallel connection of N capacitors, they all share the same voltage across them. The equivalent capacitance for such a configuration is given by

Equation1

It is crucial to note that the equivalent capacitance of N capacitors connected in parallel equals the total of their individual capacitances.

Transitioning now to the scenario where N capacitors are interconnected in a series, it is observed that the same current i, and consequently the same charge, flows through all the capacitors. The equivalent capacitance of this setup is represented as

Equation2

In contrast to the parallel configuration, the equivalent capacitance of capacitors connected in series is calculated as the reciprocal of the sum of the reciprocals of each capacitor's individual capacitance.

 Core: Electrical Engineering

Graphical and Analytic Representation of Sinusoids

JoVE 15105

Analyzing two sinusoidal voltages with equal amplitude and period but different phases on an oscilloscope, an instrument used to display and analyze waveforms, involves a three-step process.

The first step is measuring the peak-to-peak value, which is twice the amplitude of the sinusoid. This provides information about the maximum voltage swing of the waveform.

Secondly, the period and angular frequency are determined. The period is the time taken for one complete cycle of the waveform, while the angular frequency (often denoted by the Greek letter omega, ω) is the rate at which the waveform oscillates.

The third step involves measuring the voltage values for both sinusoids at a fixed point in time. This is used to determine the phase angle, which depends on whether the sinusoid has a positive or negative slope at that point.

A graphical representation can be employed to compare the two sinusoids. In this representation, the horizontal axis represents the magnitude of cosine, and the vertical axis represents the magnitude of sine. This arrangement is reminiscent of polar coordinates, where angles measured counterclockwise from the horizontal axis are considered positive, while those measured clockwise are deemed negative.

Interestingly, subtracting 90 degrees from the argument of a cosine function yields the sine function. This property can be leveraged when dealing with sinusoids.

To add two sinusoids of the same frequency - one in sine form and the other in cosine form - the graphical representation comes into play again. The hypotenuse of the right-angled triangle formed by these sinusoids represents the resultant sinusoid. The argument of this resultant sinusoid, measured from the horizontal axis, equals the cosine inverse of the ratio of the base to the hypotenuse.

In conclusion, the use of an oscilloscope coupled with a firm understanding of trigonometric principles provides a powerful toolset for analyzing sinusoidal voltages. This knowledge not only aids in understanding the behavior of AC circuits but also finds applications in numerous fields like telecommunications, signal processing, and power systems.

 Core: Electrical Engineering

Design Example: Frog Muscle Response

JoVE 15174

A student is tasked to work on an intriguing experiment involving an RL (Resistor-Inductor) circuit to study the muscle response of a frog's leg to electrical stimulation. The RL circuit plays a crucial role in this experiment, providing the means to control and measure the electrical impulses that trigger muscle contraction.

When the switch connecting the RL circuit is closed, a brief muscle contraction is observed. This is because, at a steady state, the inductor acts like a short circuit, bypassing the resistor and causing a mild, brief contraction in the frog's leg. Applying Ohm's law to the 60-ohm resistor, the student can calculate the initial current passing through the circuit.

Upon opening the switch, the RL circuit becomes source-free. Now, the current flows through the resistor, which the student models as the frog's leg. This change causes sustained muscle activity lasting for ten seconds.

The student assumes that a current of 20 milliamperes is responsible for inducing this sustained muscle response. To verify this assumption, they use the time constant of the RL circuit, which equals the ratio of inductance to resistance. The current passing through the inductor can be expressed using the calculated initial current and this time constant.

Substituting the current and time values corresponding to the sustained muscle activity into the current equation, the student can solve for the unknown resistance. This resistance value represents the modeled resistance of the frog's leg.

This experiment serves as a practical demonstration of how RL circuits can be used to study physiological responses. It illustrates the fundamental principles governing these circuits, such as transient and steady-state responses and the concept of the time constant. By drawing parallels between electrical circuits and biological systems, it also highlights the interdisciplinary nature of science, bridging the gap between physics and biology.

 Core: Electrical Engineering

Phasor Arithmetics

JoVE 15568

Phasors and their corresponding sinusoids are interrelated, offering unique insights into the behavior of alternating current (AC) circuits. One way to understand this relationship is through the operations of differentiation and integration in both the time and phasor domains.

When the derivative of a sinusoid is taken in the time domain, it transforms into its corresponding phasor multiplied by j-omega (jω) in the phasor domain, where j is the imaginary unit, and ω is the angular frequency. Conversely, when a sinusoid is integrated in the time domain, it translates into its corresponding phasor divided by j-omega in the phasor domain. These transformations provide a means to find steady-state solutions for the sinusoid without knowing the initial variable values.

Next, consider two phasors, each represented in rectangular and polar forms. To add or subtract these two phasors, their rectangular forms are used (which express the phasor as a complex number with real and imaginary parts). The real part of the resultant phasor is the sum (for addition) or difference (for subtraction) of the real parts of the two original phasors, and its imaginary part is the sum or difference of the imaginary parts of the individual phasors.

When multiplying or dividing any two phasors, their polar forms are used (expressing the phasor as a magnitude and an angle). The magnitude of the resultant phasor is the product (for multiplication) or quotient (for division) of the magnitudes of the two original phasors, and the angle of the resultant phasor is the sum or difference of the angles of the individual phasors.

Lastly, the complex conjugate of a phasor - which is obtained by changing the sign of its imaginary part - can be expressed in both rectangular and polar forms. This operation is crucial in many applications, including the computation of power in AC circuits.

In conclusion, phasors serve as a powerful mathematical tool in the study of AC circuits, simplifying analysis and solving problems that would be significantly more complex in the time domain.

 Core: Electrical Engineering

Zebrafish Microinjection Techniques

JoVE 5130

One of the major advantages to working with zebrafish (Danio rerio) is that their genetics can be easily manipulated by microinjection of early stage embryos. Using this technique, solutions containing genetic material or knockdown constructs are delivered into the blastomeres: the embryonic cells sitting atop the yolk of the newly fertilized egg. Delivery into the cytoplasm is achieved either through direct injection into the blastomere, or via natural cytoplasmic movements that occur after a solution is injected into the yolk. Successful genetic manipulations are usually followed by quantification of embryonic phenotypes in order to elucidate the genetic mechanisms of development.

This video will provide an introduction to carrying out microinjections in zebrafish embryos. The discussion begins with a review of the essential tools for the technique, including the injection apparatus and the microinjector, which controls fluid movement with pressure pulses of air. Next, important preparatory steps are demonstrated, such as the pouring of agar plates to stabilize embryos during injection and calibration of the microinjection apparatus. The injection procedure is then presented along with tips on when and where injections should be performed. Finally, applications of the microinjection technique are discussed, including gene overexpression via mRNA injection, gene silencing by delivery of antisense morpholino oligonucleotides, and the generation of transgenic zebrafish using specially engineered plasmid DNA.

 Biology II

An Introduction to Neurophysiology

JoVE 5201

Neurophysiology is broadly defined as the study of nervous system function. In this field, scientists investigate the central and peripheral nervous systems at the level of whole organs, cellular networks, single cells, or even subcellular compartments. A unifying feature of this wide-ranging discipline is an interest in the mechanisms that lead to the generation and propagation of electrical impulses within and between neurons. This subject is important not only for our understanding of the fascinating processes driving human thought, but also for our ability to diagnose and treat disorders related to nervous system malfunction.

This video will provide an introduction to the field of neurophysiology, beginning with a brief history of neurophysiological research that showcases landmark studies like Galvani’s observations of twitching frog legs and Eccles’s discovery of the chemical synapse. Next, key questions asked by neurophysiologists are introduced, followed by an overview of some prominent experimental tools used to answer those questions. The methods presented range from techniques used to investigate single cells, like patch clamping, to those that can measure activity across large regions of the brain, like electroencephalography (EEG). Finally, applications of neurophysiological research are discussed, including the development of brain-machine interfaces that allow for device control using thought.

 Neuroscience

An Overview of Gene Expression

JoVE 5546

Gene expression is the complex process where a cell uses its genetic information to make functional products. This process is regulated at multiple stages, and any misregulation could lead to diseases such as cancer.

This video highlights important historical discoveries relating to gene expression, including the understanding of how distinct combinations of DNA bases encode the amino acids that make up proteins. Key questions in the field of gene expression research are explored, followed by a discussion of several techniques used to measure gene expression and investigate its regulation. Finally, we look at how scientists are currently using these techniques to study gene expression.

 Genetics

Cell Cycle Analysis

JoVE 5641

Cell cycle refers to the set of events through which a cell grows, replicates its genome, and ultimately divides into two daughter cells through the process of mitosis. Because the amount of DNA in a cell shows characteristic changes throughout the cycle, techniques known as cell cycle analysis can be used to separate a population of cells according to the different phases of cell cycle they are in, based on their varying DNA content.

This video will cover the principles behind cell cycle analysis via DNA-staining. We will review a generalized protocol for performing this staining using bromodeoxyuridine (BrdU, a thymidine analog that is incorporated into newly synthesized DNA strands) and propidium iodide (PI, a DNA dye that stains all DNA), followed by analysis of the stained cells with flow cytometry. During flow cytometry, a single cell suspension of fluorescently labeled cells is passed through an instrument with a laser beam and the fluorescence of each cell is read. We will then discuss how to interpret data from flow cytometric scatter plots, and finally, look at a few applications of this technique.

 Cell Biology

Chromatography-based Biomolecule Purification Methods

JoVE 5683

In biochemistry, chromatography-based purification methods are employed to isolate compounds from a complex mixture. Two such methods used commonly by biochemists are size-exclusion chromatography and affinity chromatography. In size-exclusion chromatography, a column packed with porous beads separates components of a mixture based on size. On the other hand, affinity chromatography allows for a more specific separation of biomolecules by using a column that is composed of stationary phase, which contains target-specific ligands.

This video serves as an introduction to size-exclusion and affinity chromatography, as well as the concepts that govern them. A step-by-step procedure for the purification of a histidine-tagged protein by immobilized metal affinity chromatography is described. Applications for both of these chromatographic methods in biochemistry and biomedical research are also profiled.

"Chromatography" refers to a wide range of methods used to isolate a component from a complex mixture, an essential step before a biomolecule's properties and activities can be determined. Each chromatographic technique has a different mechanism for separation, depending on the sample matrix and target compound. This video will focus on the principles and operation of two methods common to biochemistry: size-exclusion and affinity chromatography.

Size-exclusion chromatography or SEC is based on the size of the compounds in the sample. A mobile phase containing the sample is added to a column with a porous material: the stationary phase. The molecules in the sample fall into 1 of 3 categories.

Molecules too large to enter the pores travel the shortest distance through the column. Any species with a molecular weight above this "exclusion limit" will exit the column at the same time. Molecules small enough to freely enter the pores will be retained the longest, and will exit the column together. The molecular weight allowing complete pore entry is the "permeation limit".

Only molecules between these limits will be separated from one another, as they spend varying amounts of time diffusing into and out of the pores. Smaller molecules are retained longer on the column because they spend more time in the stationary phase, whereas larger molecules within these limits exit earlier.

Molecular weights across 1 to 2 orders of magnitude fall within these limits. Columns are chosen with this in mind, or multiple columns can be used in series if there is a wide range of desired compounds.

Now that you've seen the theory of SEC, let's look at how it is carried out.

To begin the SEC procedure, the column must be equilibrated with deionized water and chromatography buffer. Once prepared, buffer containing the sample is injected onto the column. The buffer is then pushed through at a low flow rate. A detector monitors what exits the column to determine the presence of the desired analyte. Large molecules with molecular weights above the exclusion limit exit the column at the same time. Small fractions from the column are collected in tubes. Each fraction is tested for the quality of the target molecule by gel electrophoresis or other analytical techniques.

Now let's have a look at affinity chromatography or AC, one of the most efficient ways to purify proteins. Many biomolecules bind selectively to certain ligands-a property which AC utilizes by adhering a target-specific ligand to the stationary phase.

When the mixture flows through the column, the target molecules attach to the ligand, and the rest flow through. After the mixture has passed through the column, the target molecule can be collected through one of two elution methods based on specificity.

Biospecific elution can be said to a have a "normal-" or "reverse-role". In normal-role biospecific elution, an agent is added that competes with the adhered ligand to bind with the target biomolecule.

In reverse-role biospecific elution, an agent competes with the target to bind to the adhered ligand. The second elution type, nonspecific elution, lowers the target-to-ligand binding by changing the solution's pH, ionic strength, or polarity. If a protein does not bind to a ligand that can be immobilized, the protein can be expressed containing a "tag": short peptide sequences engineered to bind to the ligand.

One variety is immobilized metal ion affinity chromatography, "IMAC" for short, where an adhered metal ligand, like nickel or cobalt, binds to histidine residues on the modified protein. Through molecular biology techniques, target proteins are generated with repeating histidine residues, called a polyhistidine-tag, which binds to the metal via the imidazole side chain on histidine. Once bound, the protein can be collected with free imidazole via reverse-role specific elution and later used in a wide variety of downstream applications. Now that you've seen the theory of affinity chromatography, let's look at an IMAC procedure in the laboratory.

In IMAC, the stationary phase can be added directly to the mixture of mobile phase and sample, allowing the binding of the his-tagged protein. This slurry is then poured into the column, where the non-bound compounds drip into waste, while the slurry remains. The slurry container is rinsed to collect residual resin and sample, which is added to the column.

The resin is stirred to ensure the unbound components are free-flowing. Added wash buffer helps flush them away. Once all of the unbound components have been removed, the waste is replaced with a container to collect the target protein.

Buffer containing imidazole is added, stirred, and allowed to rest to unbind the target molecule. The imidazole binds to the metal, replacing and releasing the tagged protein. The freed protein is collected, and the imidazole step is repeated to ensure total collection. To further purify the sample, SEC can be run on the sample prior to analysis.

Now that we've seen the theory and procedure of these two techniques, let's look at some of the ways they're applied in the biochemical field.

A common reason to purify proteins is to study their role in disease. Cystic fibrosis is caused by defects in the cystic fibrosis transmembrane conductance regulator protein, or CFTR. After growing the protein with a his-tag in yeast, both affinity and size-exclusion chromatography allow the isolation of the protein, followed by the study of its function.

In some instances, the presence of a polyhistidine tag can change a protein's structure, thereby affecting its function. Another common tag is maltose-binding protein or MBP, which will bind to bound amylose in a column. Maltose is then used to release the complex. The MBP can then be cleaved and removed with SEC to produce the pure desired protein.

You've just watched JoVE's video on size-exclusion and affinity chromatography. It covered the theory of the techniques, went over general procedures, and covered some of the uses of the techniques.

Thanks for watching!

 Biochemistry

Fractional Distillation

JoVE 5700

Source: Laboratory of Dr. Nicholas Leadbeater — University of Connecticut 

Distillation is perhaps the most common laboratory technique employed by chemists for the purification of organic liquids. Compounds in a mixture with different boiling points separate into individual components when the mixture is carefully distilled. The two main types of distillation are "simple distillation" and "fractional distillation", and both are widely used in organic chemistry laboratories.

Simple distillation is used when the liquid is (a) relatively pure (containing no more than 10% liquid contaminants), (b) has a non-volatile component, such as a solid contaminant, or (c) is mixed with another liquid with a boiling point that differs by at least 25 °C. Fractional distillation is used when separating mixtures of liquids whose boiling points are more similar (separated by less than 25 °C).

This video will detail the fractional distillation of a mixture of two common organic solvents, cyclohexane and toluene.

 Organic Chemistry

Whole Organ Tissue Culture

JoVE 5799

Whole organs can be cultured ex vivo using specialized bioreactors with the goal of repairing or replacing entire organs. This method uses a donor organ that is stripped of all cells, leaving behind the three-dimensional structure, which is then repopulated with new cells. This video demonstrates the whole organ culture of lungs and shows how a dynamic culture that mimics the mechanical stimulation in the body is needed to induce native tissue properties.

 Bioengineering

Stereoisomers

JoVE 11721

On the basis of mirror symmetry, stereoisomers of an organic molecule can be further classified into diastereomers and enantiomers. Diastereomers are stereoisomers that are not mirror images of each other. Substituted alkenes, such as the cis and trans isomers of 2-butene, are diastereomers, as these molecules exhibit different spatial orientations of their constituent atoms, are not mirror images of each other, and do not interconvert. Here, the interconversion is suppressed due to restricted rotation around the π bond. Another class of diastereomers — those without π bonds — are molecules that have non-superposable spatial orientations and exhibit different configurations of their substituent groups at some, but not all, stereocenters. For example, cis-1,2-dimethylcyclohexane and trans-1,2-dimethylcyclohexane are diastereomers, as these molecules are not superposable and have a different configuration of the methyl and hydrogen groups at only one out of their two stereocenters.

Enantiomers are stereoisomers that are mirror images of each other. As only chiral molecules can have non-superposable mirror images, enantiomers are chiral molecules. For example, the chiral molecule 2-butanol and its mirror image are enantiomers, as these molecules exhibit non-superposable spatial orientations of their constituent atoms and are mirror images of each other. A chiral molecule and its mirror image are collectively referred to as an enantiomeric pair, or a pair of enantiomers. The enantiomer of a chiral molecule can be drawn by taking its mirror image from any position or exchanging the positions of two of the substituents at each stereocenter of the molecule.

 Core: Organic Chemistry

Isomerism in Alkenes

JoVE 11767

Alkenes like 1-butene and 2-butene exhibit constitutional isomerism, as they differ in the position of the double bond. Further, 2-butene exhibits stereoisomerism and exists as two distinct compounds differing in spatial arrangement.

An isomer is called cis-2-butene when the methyl groups are on the same side of the double bond, and the other stereoisomer, in which methyl groups are on the opposite side of the double bond, is called trans-2-butene. The cis and trans stereoisomers are not interconvertible at room temperature because of the restricted rotation of the double bond.

Figure1left Figure1right
cis-2-Butene trans-2-Butene

For tri- or tetrasubstituted compounds, the E,Z nomenclature has been adopted over the cis and trans nomenclature. The E,Z nomenclature is based on the sequence rules. The carbon atoms across the double bond are envisioned separately, and the substituent with connecting atoms having higher atomic numbers is assigned the highest priority. If two substituents have the same atomic number, the first point of difference is decisive. When the high-priority substituents on both the carbon atoms of the double bond are on the same side, the alkene has a Z configuration, whereas when the high-priority substituents are on the opposite side of the double bond, the alkene is said to have an E configuration.

Figure2left Figure2right
Z Configuration E Configuration

 Core: Organic Chemistry

Reduction of Alkenes: Asymmetric Catalytic Hydrogenation

JoVE 11786

Catalytic hydrogenation of alkenes is a transition-metal catalyzed reduction of the double bond using molecular hydrogen to give alkanes. The mode of hydrogen addition follows syn stereochemistry.

The metal catalyst used can be either heterogeneous or homogeneous. When hydrogenation of an alkene generates a chiral center, a pair of enantiomeric products is expected to form. However, an enantiomeric excess of one of the products can be facilitated using an enantioselective reaction or an asymmetric hydrogenation process using chiral homogeneous catalysts. The chiral catalysts are designed such that the metal coordinates to a chiral ligand. The most frequently used chiral ligand is BINAP [(2,2'-bis(diphenylphosphino)-1,1'-binaphthyl] — a chelating diphosphine. The metal coordinates to the two phosphorus atoms of BINAP, creating a chiral environment for itself. Such chiral catalysts have tremendous applications in pharmaceutical industries, such as the asymmetric synthesis of (S)-naproxen, an anti-inflammatory drug molecule, and the synthesis of L-dopa, a drug used to treat patients with Parkinson's disease.

Asymmetric hydrogenation is specific to the type of double bond undergoing reduction. The presence of a functional group directly adjacent to the target double bond is essential for the hydrogenation process as it aids with effective coordination of the metal.

 Core: Organic Chemistry

Microtubules

JoVE 11906

Microtubules are the thickest cytoskeletal filaments with a diameter of 25 nm. In prokaryotic organisms, microtubules are commonly found in locomotory appendages like cilia and flagella. In eukaryotic cells, microtubules form specialized extensions for moving fluid over the surface, like those found in cells lining the intestine.

Microtubules have two structurally similar globular protein subunits: α and β tubulins. In the cytosol, the α and β tubulins form a heterodimer. These αβ-heterodimers enter the centrioles, a type of microtubule organizing center (MTOCs), and assemble into a filament with the help of the γ-tubulin ring complex. Each microtubule has 13 protofilaments, each with alternating α and β tubulins.

Microtubules are polar structures with plus (+) and minus (−) ends. The end originating from the MTOCs is the minus end, while the outward-facing end is the plus end. Microtubules are dynamic structures that undergo repeated polymerization and depolymerization. However, these filaments are not flexible; they cannot flex or bend when force is applied, and they will break apart if a deforming force is sufficient.

In microtubules, both α-tubulin and β-tubulin are associated with GTP. The GTP-bound-β-tubulins at the plus end are necessary for microtubule polymerization.  GTP-α-tubulin cannot be hydrolyzed to GDP, however, GTP-β-tubulin may be. The hydrolysis of GTP to GDP breaks the lateral interaction between the protofilaments, allowing tubulin monomers to dissociate.

Microtubules have multiple roles within the eukaryotic cell, and they act as tracks for transporting cargo and vesicles within the cell due to their polar nature. Microtubule-associated motor proteins, such as kinesin and cytoplasmic dyneins, transport cargoes throughout the cell.

 Core: Cell Biology

Alcohols from Carbonyl Compounds: Grignard Reaction

JoVE 11926

Grignard reagents are one of the most commonly used reagents used to synthesize alcohols from carbonyl compounds. Grignard reagents are organomagnesium halides with a highly polar carbon–magnesium bond. Due to the partial ionic nature of the C–Mg bond, the carbon functions as a strong nucleophile and attacks electrophiles like carbonyl carbon.

Magnesium from the reagent coordinates with carbonyl oxygen, further reducing the carbonyl carbon's electron density. Thus, the carbonyl carbon is a stronger electrophile. The carbanionic group of the reagent then attacks this electrophilic carbon to form an alkoxide ion. In the following step, the alkoxide ion is protonated using dilute acid or water to give alcohol.

Different classes of alcohols are formed depending on the type of carbonyl compounds undergoing the reaction. When reacted with a Grignard reagent, formaldehyde converts to a primary alcohol, while all other aldehydes form secondary alcohols. Ketones, on the other hand, give tertiary alcohols. The carbonyl group in a carboxylic acid does not welcome nucleophilic additions with the Grignard reagent. Instead, the reagent acts as a strong base and abstracts a proton from the acid's –COOH group. The derivatives of carboxylic acid, such as esters and acid chlorides, react with two equivalents of Grignard reagent to form tertiary alcohols with a ketone intermediate.

 Core: Organic Chemistry

Satellite Stem Cells and Muscular Dystrophy

JoVE 12523

Satellite stem cells or myosatellite cells are quiescent stem cells that Alexander Mauro first identified in 1961. These cells are located between the sarcolemma, the plasma membrane of muscle fibers, and the basal lamina, the connective tissue sheath covering it. These mononucleated cells are activated in response to muscle injury, can transform into myoblasts, and may form or repair muscle fibers. Myosatellite cells can provide additional myonuclei for muscle regeneration or return to a quiescent state. Thus, myosatellite cells are a  source of muscle progenitor cells in postnatal and adult life.

After muscle injury in a normal patient, myosatellite cells can repair damaged muscle cells; however, in patients with muscular dystrophy, the satellite cells fail to repair the damaged muscle.  The cell will undergo apoptosis, further causing muscle wastage and eventually death. Muscular dystrophy is a group of diseases, with Duchenne muscular dystrophy (DMD) being the most common. DMD is caused by a mutation in the gene encoding dystrophin. The dystrophin protein connects the cytoskeleton of the muscle cell with the extracellular matrix and stabilizes the plasma membrane. In DMD patients, the dysfunctional dystrophin makes the plasma membrane fragile. Excess calcium ions enter the cell, damaging the muscle fibers. Additionally, white blood cells, such as macrophages, damage the connective tissue surrounding the muscle cell.

Until recently, no treatment was available for muscular dystrophy; however, gene therapy appears to be a promising solution in clinical research trials. The delivered gene encodes for microdystrophin, a shorter version of dystrophin that can function to restore normal muscle activity.

 Core: Cell Biology

Actin Treadmilling

JoVE 12599

Actin filaments undergo polymerization and depolymerization from either end. The polymerization and depolymerization rates depend on the cytosolic concentration of free G-actins. The polymerization rate is generally higher at the plus or barbed end, while the depolymerization rate is higher at the minus or pointed end. At a steady state, critical concentration describes the concentration of free G-actin monomers at which the polymerization rate at the plus end is equal to that of the depolymerization rate from the minus ends. Actin treadmilling is the continuous addition at one end and dissociation from the other end of the G-actin monomers. 

Actin treadmilling is crucial to several functions in eukaryotic cells, such as cell migration, endocytosis, and exocytosis. Treadmilling occurs continuously, even in resting cells, and is responsible for up to 50% of energy consumption in most cells.

Actin treadmilling depends primarily on three factors,  the ATP hydrolysis rate within the F-actin, the polarity in the polymerization and depolymerization rate at both plus and minus ends, and the cytosolic concentration of free G-actins. ATP-G-Actins bind faster at the plus end and undergo gradual ATP hydrolysis, forming intermittent ADP-Pi-Actins and, eventually, ADP-Actins. ADP-G-actins are loosely associated with the monomers with the F-actins and readily dissociate. At the plus end, the polymerization rate is higher than depolymerization, while at the minus end, the depolymerization rate of G-actin is higher than their polymerization rate. The cytosolic G-actin concentration during the actin treadmilling remains intermediate between that of the critical concentrations of plus and minus ends of F-actin.

 Core: Cell Biology

Amplifying Signals via Enzymatic Cascade

JoVE 13162

When a ligand binds to a cell-surface receptor, the receptor's intracellular domain changes shape, which may either activate its enzyme function or allow its binding to other molecules. The initial signal is amplified by most signal transduction pathways. This means that a single ligand molecule can activate multiple molecules of a downstream target. Proteins that relay a signal are most commonly phosphorylated at one or more sites, activating or inactivating the protein. Kinases catalyze the transfer of phosphate groups to the proteins and are specific for different target proteins.

Signal termination is essential for regulating the enzymatic cascades, and aberrant signal termination is often seen in tumor cells. Inactivation of phosphorylated signaling proteins is carried out by protein phosphatases that dephosphorylate the proteins. Therefore, phosphatases are considered the key regulators of signal transduction pathways. For example, several tyrosine phosphatases are recruited to the membrane when a ligand binding to the receptor stimulates receptor phosphorylation. SHP-1, a tyrosine phosphatase with an SH2 domain, binds to phosphotyrosines on activated cytokine receptors. On activation by JAK2- dependent phosphorylation, SHP-1 dephosphorylates specific JAKs and STATs to turn off the JAK/STAT signaling pathway.

In the case of G protein signaling, G protein itself possesses a GTPase activity that hydrolyzes its bound GTP into GDP and turns off the cascade. β-adrenergic receptor kinase, together with β-arrestin, also turns off G protein signaling. The receptor phosphorylated by β-adrenergic receptor kinase binds to β-arrestin, blocking the receptor from interacting with G proteins. Lastly, phosphodiesterases cause a reduction in the levels of second messenger cAMP, terminating the signaling through G protein-coupled receptor.

 Core: Cell Biology

G Protein-coupled Receptors

JoVE 13320

G Protein-Coupled Receptors or GPCRs are membrane-bound receptors that transiently associate with heterotrimeric G proteins and induce an appropriate response to sensory stimuli such as light, odors, hormones, cytokines, or neurotransmitters.

GPCRs are also called heptahelical, 7TM, or serpentine receptors, and consist of seven (H1-H7) transmembrane alpha-helices that span the bilayer to form a cylindrical core. The transmembrane helices are connected by three extracellular loops and three cytosolic loops. Together with the extracellular loops, the transmembrane alpha-helices form the central ligand-binding pocket of GPCR. In contrast, the third cytosolic loop functions as the heterotrimeric G protein binding site.

Ligand binding activates the GPCRs as it undergoes a conformational change and also binds heterotrimeric G proteins with high affinity. An activated GPCR can  bind and activate multiple G proteins to amplify the signal. G proteins, in turn, bind and activate downstream effectors and bring about a cellular response.

Although structurally, all mammalian GPCRs consist of seven transmembrane alpha-helical domains, they differ considerably in their sequence and functionality. GPCRs are broadly categorized into five classes, including Class A (rhodopsin-like), Class B (secretin receptor-like or B1), Class B2/ adhesion type,  Class C (glutamate receptor-like), and Class F (frizzled-like).

  • Class A forms the largest subfamily of GPCRs that includes rhodopsins and beta-adrenergic receptors.
  • Class B comprises the hormone-binding receptors like glucagon, parathyroid hormone, and vasoactive intestinal peptide (VIP) receptors.
    • The adhesion or B2 receptor class includes the adhesion G protein-coupled receptors or ADGR groups of receptors such as ADGRL1 and ADGRG1 that are essential for cell adhesion and migration.
  • Class C includes calcium-sensing receptors, gamma-aminobutyric acid (GABA) type B receptors, metabotropic glutamate receptors, and several taste receptors. Unlike the others, this group uses a characteristic venus fly trap module for ligand binding.
  • Class F, also called frizzled-like, includes smoothened or Smo receptors and functions in embryonic development.

Overall, humans consist of more than 800 GPCRs. Many of these detect hormones, growth factors, or endogenous ligands, while several others are involved in olfactory and gustatory responses.

Thus, GPCRs regulate critical physiological pathways and are an excellent drug target for treating diseases like diabetes, cancer, obesity, depression, or Alzheimer's. Nearly 35% of approved drugs implement their therapeutic effects by selectively interacting with specific GPCRs. One commonly used class of drugs, beta-blockers,  target beta-adrenergic receptors and treat conditions like hypertension, cardiac arrhythmia, and anxiety. GPCRs provide an effective target to create an arsenal for a varied range of diseased conditions.

 Core: Cell Biology

Cellulose and Pectic Polysaccharides

JoVE 13367

 Every plant cell has a cell wall that protects the cell, provides structural support, and gives the cell shape. Cellulose, the main structural component of the plant cell wall, makes up over 30% of plant matter. It is the most abundant organic compound on earth.  Cellulose is an unbranched polysaccharide composed of linear chains of glucose molecules linked by β (1→4) glycosidic bonds.

As a cell matures, its cell wall specializes according to its type. For example, the parenchyma cells of leaves possess only a thin, primary cell wall. Collenchyma and sclerenchyma cells, on the other hand, mainly occur in the outer layers of a plant's stems and leaves. These cells provide strength and support by partially thickening their primary cell wall or depositing a secondary cell wall. Some plants, such as trees and grasses, deposit a secondary cell wall around mature cells. Secondary cell walls typically contain three distinct layers: in each layer, the cellulose microfibrils are organized in different orientations.

The cell wall contains other polysaccharides, such as pectin. Pectic polysaccharides are heteropolysaccharides that predominantly have a galacturonic acid sugar backbone with acid sugar and neutral sugar side chains. During cell growth and expansion, pectic polysaccharides are deposited in the middle lamella, followed by the primary and secondary cell walls. They are crucial in wall hydration and cell-cell adhesion. They also influence wall porosity, regulate ion transport, and help in the morphogenesis of plants by impacting the enzymes in the cell wall and their water-holding capacity.

Homogalacturonan, the most prominent type of pectin, accounts for around 60% of total pectin in cell walls.

Adapted from openstax biology 2e

 Core: Cell Biology

MALDI-TOF Mass Spectrometry

JoVE 13384

Mass spectrometry is a powerful characterization technique that can identify and separate a wide variety of compounds ranging from chemical to biological entities, based on their mass-to-charge ratio (m/z). The instruments that allow this detection, known as mass spectrometers, have three components: an ion source, a mass analyzer, and a detector. These spectrometers differ based on the nature of their ion source and analyzers.

Matrix-assisted laser desorption ionization (MALDI) is a commonly used mass spectrometer. It enables the soft ionization of biological molecules like peptides, lipids, and saccharides without fragmenting them. This means they can be analyzed intact as the sample integrity is maintained.

In MALDI, the sample is mixed with a compatible matrix that is organic molecules. When the laser strikes the sample, the matrix functions as a mediator for energy absorption, yielding the intact sample analyte ions. The matrix molecules energetically ablate from the sample surface, absorb the laser energy and carry the sample molecules into the gas phase. The sample molecules are usually ionized during the ablation process to carry a single positive charge.

The desorbed sample ions carrying the charge are then directed towards the mass analyzer, usually a time-of-flight (TOF). The ions pass through a field-free drift region which allows separation based on size, allowing the smaller ions to reach the detector first.

MALDI-TOF instruments are often connected with a reflectron that reflects ions increasing the ion flight path. This increased flight time between ions of different m/z allows better sample resolution where ions of the same mass reach the detector simultaneously. MALDI sources can be coupled with other analyzers, including triple quadrupoles and Fourier transform ion cyclotron resonance (FT-ICR) mass spectrometers. The MALDI-FT-ICR is known for its high mass resolution, wherein the sample ions carrying the charge move to an ICR cell and cyclotron in a magnetic field while separating on m/z.

With a couple of limitations, like, the possibility of photodegradation by the laser for some sample types, absorption of the laser by a few fluorescent analytes that could affect the results, an acidic matrix that could alter sample nature, etc., overall, this technique has vast applications.  It can analyze protein samples isolated by electrophoresis or chromatography, peptides, carbohydrates, DNA, identify microorganisms from clinical samples and ascertain biomolecules in tissues, among many others.

 Core: Cell Biology

Scanning Electron Microscopy

JoVE 13400

A scanning electron microscope (SEM) is used to study the surface features of a sample by using an electron beam that scans the sample surface in a two-dimensional manner. Typically, areas between ~1 centimeter to 5 micrometers in width can be imaged. SEM can be used to image bacteria, viruses, tissues as well as larger samples like insects. Conventional SEM gives a magnification ranging from 20X to 30,000X and spatial resolution of 50 to 100 nanometers.

Fundamental Principles

Accelerated electrons released by the electron gun have high kinetic energy (ranging from 5-30 keV). Electron-sample interactions lead to deceleration of the electrons and dissipation of the energy in the form of different signals. The electrons undergo two types of scattering: elastic and inelastic. Inelastic scattering causes the emission of secondary electrons. These low-energy electrons (~50 eV) are the outer shell electrons of the sample atoms that acquire just enough energy to leave the atom's surface. Only topographical information is provided by the scattering of secondary electrons since the energy level of the electrons leaving from the internal regions of the sample is too low to exit the sample surface. 

X-rays are also generated by inelastic collisions of the incident electrons with electrons in discrete orbitals (shells) of atoms in the sample. As the excited electrons return to lower energy states, they yield X-rays of a fixed wavelength (related to the difference in energy levels of electrons in different shells for a given element). Thus, characteristic X-rays are produced for each element in a mineral excited by the electron beam.

Elastic scattering, on the other hand, is not caused by dislodged electrons from the sample atoms. The principal beam of electrons is backscattered after interaction with the nucleus. These electrons do not change their energy or speed but change their direction based on their interaction with the nucleus. Detection of these electrons provides compositional information, and their varying contrast upon interaction with atoms of different atomic weights allows the user to distinguish differences in sample composition. In biological samples, this can be used to study embedded or attached nanoparticles and nanostructures with heavier atomic weights, such as gold or iron. 

 Core: Cell Biology

Renewal of Intestinal Stem Cells

JoVE 13464

The intestinal epithelial lining rapidly renews every 4 to 5 days. The renewal is facilitated by intestinal stem cells (ISCs) located at the base of the crypt– a gland located at the bottom of each villus. ISCs divide asymmetrically to form new stem cells and progenitor daughter cells. The daughter cells are called transit-amplifying (TA) cells which move upwards along the crypt and either differentiate into absorptive cells– the enterocytes or secretory cells– including the goblet, entero-endocrine, and Paneth cells.

The differentiated cells have distinct functions. Absorptive cells or enterocytes comprise 90% of the intestinal epithelial cell population. They are involved in the absorption and transport of nutrients and electrolytes. Goblet cells secrete mucus which protects the intestinal mucosa. The enteroendocrine cells secrete hormones that help in cell proliferation and digestive activities. Paneth cells secrete antibacterial proteins such as lysozyme and are also involved in the proliferation of stem cells.

Intestinal stem cell proliferation and migration were determined using a labeling method in which radioactive thymidine was injected into the experimental mice's gut tissues. The actively dividing cells incorporate the radioactive thymidine into their DNA during S-phase. When the same cell divides, the intensity of the label gets reduced as it is distributed between the new daughter cells, which can then be quantified. Based on this experiment, it was found that the rapidly dividing cells, also known as transit-amplifying cells, are present in the middle and upper parts of the crypt, and the slowly dividing stem cells are present in the bottom part of the crypt placed between the Paneth cells. The differentiated mature epithelial cells, except for Paneth cells, are present at the finger-like projections or villi.

Two models were proposed for the identity and location of ISCs– the stem cell zone model and the +4 cell model. Cheng and Leblond proposed the stem cell zone model in which they found that stem cells are interspersed between the Paneth cells. They named these cells ‘crypt base columnar’ or CBC cells. Chris Potten and colleagues proposed the +4 stem cell model. Using the DNA labeling method, they found that these cells are located at the 4th position from the Paneth cells, hence the name ‘+4 cells’.

Recent studies have shown that the CBC cells are the ISCs that participate in renewing intestinal epithelium, whereas the +4 cells are quiescent and become activated upon injury or stress. In such cases, the +4 cells give rise to progenitor cells that eventually replace the cells damaged by stress.

 Core: Cell Biology

Somatic to iPS Cell Reprogramming

JoVE 13481

Reprogramming alters the gene expression in somatic cells, transforming them into induced pluripotent stem (iPS) cells over several generations. Scientists can reprogram cells by introducing genes for four transcription factors—Oct4, Sox2, Klf4, and c-Myc (OSKM) by viral or non-viral methods. These factors are also known as Yamanaka factors after Shinya Yamanaka, who first generated iPS cells using mouse skin cells. Yamanaka was awarded the Nobel Prize in Physiology or Medicine in 2012 for this discovery.

The expression of OSKM factors brings about several cellular changes in different phases. The initiation phase downregulates genes specific to the somatic cell, upregulates genes involved in proliferation, and reactivates telomerase. Cells such as fibroblasts undergo a mesenchymal to epithelial transition, where they acquire an apical-basal polarity and express epithelial cell markers, such as cadherin, vimentins, and tight junctions. The intermediate phase involves the activation of genes required for pluripotency. Cells undergoing reprogramming use glycolysis preferentially over oxidative phosphorylation for ATP generation. This change occurs because reprogramming factors transform the elongated mitochondria into spherical ones, with very few cristae. The maturation phase induces epigenetic changes and cytoskeletal remodeling.

The entire reprogramming process alters the expression of around 1500 genes. After reprogramming, less than 1% of the cells become pluripotent. This percentage can be increased by altering the chromatin structure, repressing the expression of proteins, such as p53, that promote cell senescence, and suppressing signaling pathways or enzymes that are barriers to reprogramming.

 Core: Cell Biology

Significance Testing: Overview

JoVE 14514

Significance testing is a set of statistical methods used to test whether a claim about a parameter is valid. In analytical chemistry, significance testing is used primarily to determine whether the difference between two values comes from determinate or random errors. The effect of a particular change in the measurement protocol, analyst, or sample itself can cause a deviation from the expected result. In the case of a suspected deviation/outlier, we need to be able to confirm mathematically that the deviation comes from a determinate source and that the observation with the deviation can be logically omitted from the analysis.

Two hypotheses are used as criteria for significance testing. The null hypothesis (H0) states that the values being compared do not differ from each other significantly. In other words, if any difference exists between two values, it is ascribed to an indeterminate error. The alternate hypothesis (HA) states that the compared values are not equal, and the difference is more significant than can be explained by indeterminate error.

Before the test is performed, the hypotheses need to be stated, and a significance level (α) needs to be set. The test statistic, based on the sample mean and standard deviation, is then calculated and compared to the tabulated values, which are set at particular significance levels and defined as one- or two-tailed. If the calculated test statistic exceeds the critical values (tabulated statistic), the null hypothesis is rejected, and we state that the difference between the two values cannot be explained by random, indeterminate error.

In one-tailed significance testing, the alternative hypothesis can specify that the observed value is either higher or lower than the expected value, but not both. In two-tailed significance testing, the alternative hypothesis can simply state that the observed value is not equal to the expected value, with no regard to the direction.

Significance testing can be used on different statistical parameters of one or more data sets. Tests are given different names depending on the parameters or purpose. Significance testing is frequently applied to compare an observed value with the mean or compare two means from two different data sets. These tests are known as t-tests. Significance tests can also be performed on the variance of two data sets. In this case, the test is known as an F-test. If a significance test is used to identify outliers, the test is called a Q-test.

 Core: Analytical Chemistry

Acid–Base Equilibria: Activity-Based Definition of pH

JoVE 14530

For an ideal solution, the pH is defined as the negative logarithm of the hydrogen ion concentration. For a non-ideal solution, an accurate measurement of the pH must consider the negative logarithm of the hydrogen ion activity rather than concentration. In such a solution, the pH can be more accurately defined as the negative logarithm of a product of the hydrogen ion concentration and its activity coefficient.

In solutions of very low ionic strength—for example, pure water—the activity coefficient of the hydrogen ion is close to one when the ionic strength of the solution increases due to the addition of an electrolyte that does not donate or accept a proton. This results in a slight decrease in the pH of the solution. In other words, the addition of an electrolyte increases the hydrogen ion activity, or the effective hydrogen ion concentration in the solution, which decreases the pH of the solution.

 Core: Analytical Chemistry

Composition of Polyprotic Acid Solutions as a Function of pH

JoVE 14546

Polyprotic acids of the type H2M constitute two ionizable protons. As a result, on titration with a base, they exhibit two equivalence points in the titration curve. During titration, the species H2M, HM, and M2− will be present in the solution at different points. The fractions of H2M, HM, and M2− present at the various instances of the titration are denoted by α0, α1, and α2, respectively.

A graph with the alpha values is plotted against the volume of base added during titration. Here, a value of 0.7 for the α0 at the beginning of the titration suggests that 70% of the solution is H2M, with the remaining 30% as HM. On adding the base, the fraction of H2M, α0, decreases to nearly zero at the first equivalence point. Simultaneously, α1, representing the fraction of HM, increases and approaches unity. As more base is added, the fraction of HMdecreases and reaches zero at the second equivalence point while the fraction of M2−, represented by α2, approaches unity.

 Core: Analytical Chemistry

Precipitation Processes

JoVE 14585

The experimental conditions in a gravimetric analysis should be optimized to maximize the particle size and purity of the obtained precipitate. Ideally, the concentration of the precipitating reagent should be low with effective stirring to maintain low relative supersaturation for the growth of large crystals. In homogeneous precipitation, the precipitant is slowly generated by a chemical reaction in the solution to avoid local reagent excesses. For example, urea decomposes gradually to release hydroxide ions in the precipitation of aluminum as its hydroxide. At higher temperatures, solubility increases, and supersaturation decreases, resulting in larger crystals.

The precipitate may take up substances normally soluble in the mother liquor, resulting in coprecipitated impurities that could be adsorbed or absorbed. To minimize these impurities, crystalline precipitates are allowed to stand in the hot mother liquor in a process called digestion (Ostwald ripening). During digestion, small crystals tend to dissolve and precipitate onto the surface of larger crystals, while adjacent particles form bridges. This slow recrystallization increases particle size and decreases the effects of coprecipitation. 

 Core: Analytical Chemistry

Precipitation and Co-precipitation

JoVE 14618

Precipitation and coprecipitation methods can be used to separate a mixture of ions in a solution. In qualitative inorganic analysis, ions that form sparingly soluble precipitates with the same reagent are separated based on the differences in solubility products. For example, consider the separation of Cu(II) and Fe(II) ions by precipitation as insoluble sulfides. First, copper(II) sulfide is precipitated by the addition of acidic H2S, where the dissociation of H2S is suppressed. Adding H2S increases the product of sulfide ion and copper ion concentrations. When this product exceeds the solubility product of copper(II) sulfide, precipitation occurs. On the other hand, H2S alone cannot be used to precipitate iron(II) sulfide because its solubility product is higher than that of copper(II) sulfide. Successful precipitation of FeS requires the addition of aqueous ammonia.

Coprecipitation is a method used to remove otherwise soluble contaminants by using a substance that can coprecipitate with the contaminants. This technique can be used to isolate trace components from a solution. For instance, the tendency of arsenic to coprecipitate with ferric hydroxide is exploited in a process called gathering. During this process,  Fe(II) or solid iron particles are added to water contaminated by arsenic. The iron particles are then oxidized for several hours to obtain ferric hydroxide precipitates, which now contain arsenic. Finally, the arsenic-bearing precipitate is removed by repetitive filtration to produce safe drinking water.

 Core: Analytical Chemistry

Motor Unit Stimulation

JoVE 14847

When the neuron of a motor unit fires an action potential, it triggers a series of events, leading to a twitch contraction in the muscle fibers. The process of excitation-contraction coupling is crucial in relaying the action potential to the muscle fibers.

The latent period of contraction marks the onset of excitation-contraction coupling, when the action potential propagates across the sarcolemma, preparing the muscle fibers for contraction. As the fibers enter the contraction phase, the tension rapidly peaks, resulting in a forceful contraction. Although a single twitch contraction is detectable in an experimental setup, it is generally insufficient for most functional activities. Instead, an active muscle undergoes a continuous series of smooth contractions, known as graded muscle contractions, which vary in strength. These contractions are crucial in allowing us to perform a wide range of precise movements.

Temporal Summation

One way muscle contractions can be graded is through temporal summation by increasing the stimulation frequency. In this phenomenon, successive stimuli arrive before the complete relaxation of the muscle fiber from the previous twitch. Each subsequent impulse cumulatively adds to the residual force from the prior contraction, causing the muscle tension to build up with each successive stimulus. The result is a smoother and stronger muscle contraction compared to the response elicited by a single stimulus, allowing more precise and forceful movements.

Motor Unit Recruitment

Another way muscle contractions can be graded is through variable recruitment of motor units. The strength of a muscle contraction is determined by how many motor units are activated and the frequency of their activation. When a light force is needed, the nervous system recruits smaller motor units, which consist of fewer muscle fibers and are more easily excited. These small motor units are often sufficient for tasks requiring precision rather than strength, such as writing or gently lifting an object. As the demand for force increases, larger motor units, with more and larger fibers, are progressively recruited to increase the muscle's power output.

This ability to vary the stimulation frequency and recruit different motor units allows muscles to contract with a range of forces, providing the precise control necessary for complex tasks and the strength required for more demanding activities.

 Core: Anatomy and Physiology

Muscles of the Eye

JoVE 14867

The muscles of the eye are sophisticated structures that control eye movement and focus, allowing for the precise and rapid adjustments necessary for vision. The human eye is controlled by ten muscles — six extraocular muscles, three intraocular muscles, and one primary eyelid retractor muscle.

Extraocular Muscles

The six extraocular muscles surround the eyeball and control its movements. They are responsible for a wide range of eye motions, including looking up, down, left, right, and rotating the eye. These include the lateral rectus, which moves the eye outward, and the medial rectus, which brings it inward toward the nose. The superior rectus and inferior rectus elevate and depress the eye, respectively, while the superior oblique and inferior oblique muscles rotate the eye. The four recti — superior, inferior, lateral, and medial recti — start from a common tendinous ring that encircles the optic foramen and attaches to the sclera. However, the superior oblique muscle starts above the optic foramen, while the inferior oblique starts on the maxilla below the orbit, and both muscles connect laterally to the sclera. These muscles work in synchronized pairs, allowing for coordinated focus and enabling critical functions like depth perception.

Intraocular Muscles

The intraocular muscles are located within the eye and are responsible for internal adjustments. The ciliary muscle plays a vital role in focusing vision by controlling the shape of the lens; it contracts to thicken the lens for nearby objects and relaxes for distance viewing. The iris contains two muscles: the sphincter muscle, which constricts the pupil in bright light to reduce light entry and protect the retina, and the dilator muscle, which widens the pupil under low light conditions to allow more light to enter, enhancing vision.

Primary Eyelid Retractor

Apart from these, there is a muscle called the levator palpebrae superioris, which originates in the wing of the sphenoid bone and attaches to the upper eyelid. This muscle is responsible for opening the eye. It works against the facial muscle, orbicularis oculi, responsible for closing or protruding the eyelid.

 Core: Anatomy and Physiology

Functions of the Nervous System

JoVE 14884

The nervous system is responsible for coordinating and regulating the body's functions. It functions through three main processes: sensory, integrative, and motor processes. Sensory function involves the detection and transmission of information about internal and external stimuli from sensory receptors to the CNS. The CNS processes this information through an integrative function, where it interprets and makes decisions based on the incoming sensory information. Finally, the motor function involves the transmission of instructions from the CNS to the effector organs (muscles and glands) to produce a response.

Nervous tissue is composed of two types of cells: neurons and glial cells. Neurons are excitable cells responsible for transmitting nerve impulses, which are electrical signals that travel through the nervous system. These impulses are generated by the movement of ions across the neuron's membrane, which creates a difference in electrical charge called the resting potential. When the neuron is stimulated, this resting potential is altered, triggering the generation of an action potential. The action potential travels down the length of the neuron's axon and is propagated to the next neuron or effector organ through the release of neurotransmitters. Glial cells provide support and nourishment to neurons and help to maintain the extracellular environment necessary for proper nerve impulse transmission.

In an example where an athlete sees a ball and kicks it, light from the ball enters the eye and is focused on the retina. Photoreceptor cells in the retina (rods and cones) convert the light into electrical signals (graded potentials) that are transmitted through the optic nerve to the brain. The visual signals from the retina are sent to the visual cortex in the occipital lobe of the brain. Here, the signals are integrated with other sensory information and memories, allowing the brain to recognize the ball and decide to kick it. The motor cortex in the frontal lobe of the brain sends a signal down the motor neurons in the spinal cord to the muscles in the leg, where it causes the muscles to contract and kick the ball. In this example, sensation occurs when graded potentials are generated in the photoreceptor cells of the retina and travel through the optic nerve to the brain. In the brain, the signals are integrated, and an action potential is generated in the motor cortex, which travels down the spinal cord to the muscles, causing a response. At the synapses between neurons, the electrical signal is converted to a chemical signal via neurotransmitters that cross the synapse and trigger the generation of a new electrical signal in the next neuron.

 Core: Anatomy and Physiology

Neurogenesis and Regeneration of Nervous Tissue

JoVE 14901

In the CNS, neurogenesis, the birth of new neurons from stem cells, is limited to the hippocampus in adults. In other regions of the brain and spinal cord, neurogenesis is almost non-existent due to inhibitory influences from neuroglia, especially oligodendrocytes, and the absence of growth-stimulating cues. The myelin produced by oligodendrocytes in the CNS inhibits neuronal regeneration. Furthermore, astrocytes proliferate rapidly after neuronal damage, forming scar tissue that physically blocks regeneration. It follows that injuries to the brain or spinal cord are typically irreversible.

However, in the PNS, repair is possible if the cell body is intact and Schwann cells remain active. Post-injury, the Nissl bodies in a neuron start to disintegrate, a process known as chromatolysis. Within a few days, the region distal to the damaged axon swells and breaks into fragments, and the myelin sheath deteriorates. This process is called Wallerian degeneration. Despite these changes, the neurolemma remains intact. Macrophages clear the debris, and RNA and protein synthesis increases, promoting the rebuilding or regeneration of the axon. The Schwann cells multiply and may form a regeneration tube across the injured area, guiding the growth of a new axon. However, if the injury gap is too large or filled with collagen fibers, new axons cannot grow.

 Core: Anatomy and Physiology

Spinal Cord: Gross Anatomy

JoVE 14917

The spinal cord resides within the protective confines of the vertebral column. It is the main pathway for information traveling between the brain and the body. It plays a fundamental role in nearly all bodily functions, from simple reflexes to complex motor movements. The spinal cord begins at the medulla oblongata at the base of the brainstem and extends downward, terminating at the conus medullaris near the first and second lumbar vertebrae. The spinal cord's length in adults is approximately 45 cm (18 inches), significantly shorter than the vertebral column, which continues to extend beyond the termination of the spinal cord.

Two key enlargements along the spinal cord are where the nerve roots concentrate to innervate the limbs. These enlargements are crucial for limb movement and sensory processing. The cervical enlargement extends from the fourth cervical to the first thoracic vertebra and corresponds with the nerves controlling the upper limbs. The lumbosacral enlargement, located between the ninth thoracic and twelfth thoracic vertebrae, innervates the lower limbs.

The spinal cord concludes its downward trajectory in a cone-shaped taper known as the conus medullaris. From this point, the filum terminale, a thin strand of fibrous tissue, extends to anchor the spinal cord to the coccyx, stabilizing it within the vertebral column. Below the conus medullaris, a bundle of nerve roots known as the cauda equina ("horse's tail") fans out to reach various parts of the lower body. These nerve roots provide motor and sensory functions to the legs and lower parts of the body.

Enveloping the spinal cord are three layers of spinal meninges, which offer protection, physical stability, and shock absorption. The innermost layer, the pia mater, directly adheres to the surface of the spinal cord. The arachnoid mater, the middle layer, provides a protective barrier. It is separated from the pia mater by the subarachnoid space, which contains cerebrospinal fluid. The outermost layer, the dura mater, is a tough protective sheath that forms the outermost covering of the spinal cord.

 Core: Anatomy and Physiology

Major Somatic Sensory Pathways

JoVE 14935

Sensory impulses related to touch, pressure, vibration, and proprioception from various body parts, such as the limbs, trunk, neck, and posterior head, travel to the cerebral cortex through the posterior column-medial lemniscus pathway. The pathway’s name derives from the two white-matter tracts that convey the impulses: the spinal cord's posterior column and the brainstem's medial lemniscus. First-order sensory neurons extend their axons into the spinal cord, forming the posterior columns consisting of two tracts, the gracile fasciculus and the cuneate fasciculus. Sensations from the lower part of the body travel through the gracile fasciculus tract, while sensations from the upper parts travel through the cuneate fasciculus tract. The axons synapse with second-order neurons located in the gracile or cuneate nucleus of the medulla, respectively, which then cross to the opposite side of the medulla and enter the medial lemniscus tract. The axons of the second-order neurons synapse with third-order neurons in the thalamus, which then project their axons to the primary somatosensory cortex.

Sensory impulses for pain, temperature, itch, and tickle from body parts such as the limbs, trunk, neck, and posterior head travel up to the cerebral cortex through the anterolateral or the spinothalamic pathway. The first-order neurons of this pathway connect sensory receptors of the body parts with the spinal cord. These first-order neurons have their cell bodies located in the posterior root ganglion, and their axon terminals synapse with second-order neurons located in the posterior gray horn of the spinal cord. The axons of the second-order neurons cross to the opposite side of the spinal cord and ascend to the brainstem as the spinothalamic tract. These axons end in the ventral posterior nucleus of the thalamus, where they synapse with third-order neurons. The axons of these third-order neurons project to the primary somatosensory area on the same side of the cerebral cortex as the thalamus.

The anterior and posterior spinocerebellar tract are the primary pathways through which proprioceptive impulses reach the cerebellum. These pathways transmit sensory information related to posture, balance, and coordination of skilled movements. The neurons of this pathway either do not cross over or cross over twice, so they reach the cerebellum on the same side as that of the impulse. Although these sensory impulses are not consciously perceived, they regulate movements and maintain balance.

 Core: Anatomy and Physiology

Sympathetic Activation

JoVE 14952

The sympathetic division can influence tissues and organs by releasing norepinephrine at peripheral synapses and distributing epinephrine and norepinephrine through the bloodstream. In times of crisis or stress, sympathetic activation occurs, which is regulated by sympathetic centers in the hypothalamus. As a result, sympathetic activation prepares the body for physical exertion, rapid ATP production, and heightened alertness, allowing individuals to respond effectively to challenging or threatening situations.

Physiology of Sympathetic Activation

During sympathetic activation, various physiological changes take place. The individual experiences increased alertness due to stimulation of the reticular activating system, leading to feeling "on edge." There is a surge of energy and euphoria, often accompanied by a temporary insensitivity to pain and a disregard for danger. The cardiovascular and respiratory centers in the brainstem become more active, resulting in increased heart rate, blood pressure, breathing rate, and depth of respiration.

Muscle tone is elevated, and the person may appear tense and may even begin to shiver. The mobilization of energy reserves occurs through the release of lipids by adipose tissue and the accelerated breakdown of glycogen in muscle and liver cells. These changes, along with the peripheral alterations, prepare the individual to cope with the stressful situation.

Sympathetic stimulation has longer-lasting and more widespread effects compared to parasympathetic stimulation. This effect is due to the extensive divergence of sympathetic postganglionic axons, allowing for simultaneous activation of multiple tissues. Additionally, norepinephrine lingers in the synaptic cleft longer than acetylcholine, and the release of epinephrine and norepinephrine from the adrenal medullae intensifies and prolongs the sympathetic responses.

 Core: Anatomy and Physiology

Equilibrium and Balance

JoVE 14971

The inner ear assumes dual functionalities of auditory perception and equilibrium maintenance. The vestibule is the organ responsible for balance. This organ contains mechanoreceptors, specifically hair cells, endowed with stereocilia, which aid in deciphering information regarding the position and motion of our heads. Two intrinsic components, the utricle and saccule, help perceive head position, while the semicircular canals track head movement. Neurological messages initiated in the vestibular ganglion are sent to the brainstem and cerebellum via the vestibulocochlear nerve.

The macula, a tissue found in the utricle and saccule, consists of surrounding support cells encircling the hair cells. The stereocilia, extensions of the hair cells, are embedded in a gelatinous substance called the otolithic membrane. This membrane is topped with a layer of calcium carbonate crystals or otoliths. The otoliths render the otolithic membrane heavy, causing it to move independently from the macula during head movements. In the event of a head tilt, the otolithic membrane, which slides over the macula, is influenced by gravity. This movement subsequently results in the bending of the stereocilia, causing specific hair cells to depolarize while others hyperpolarize. The brain decodes the precise head position through the depolarization pattern of hair cells.

The vestibule extends to form three ring-like structures, the semicircular canals. These are arranged in different planes: one horizontally and the other two vertically at approximately 45 degrees relative to the sagittal plane. Each canal base is connected to a dilated region called the ampulla, which houses hair cells that respond to rotational head movements. The stereocilia of these hair cells extend into a structure called the cupula, located on top of the ampulla. When the head rotates in a plane corresponding to a semicircular canal, the fluid inside the canal lags, causing the cupula to deflect oppositely to the head movement. The semicircular canals encompass several ampullae arranged both horizontally and vertically. This arrangement allows the vestibular system to decipher the direction of diverse head movements within three-dimensional (3-D) space.

Common disorders of the vestibular system:

Motion sickness:

The fundamental cause of motion sickness resides in the complex interactions between our sensory systems. The vestibular system in the inner ear plays a crucial role in maintaining balance and spatial orientation. During travel, the visual input may not align with the signals from the vestibular system, causing a sensory conflict that leads to the symptoms of motion sickness. This disparity between visual and vestibular inputs triggers neurochemical changes within the brain. Specifically, the neurotransmitter histamine is released in higher amounts, contributing to the onset of motion sickness symptoms.

Several factors can influence the onset of motion sickness. Vehicle type is one such factor; some people may be more susceptible to car motion sickness, while others may experience it more frequently on boats or airplanes. The route of travel can also contribute to the severity of motion sickness. For example, winding roads or turbulent flights may trigger more intense symptoms. Individual susceptibility is another key factor. Some people are naturally more prone to motion sickness than others, and various genetic and environmental factors can influence this.

Symptoms of motion sickness can include nausea, vomiting, dizziness, and headaches. These symptoms can significantly impact an individual's daily life, particularly if they need to travel frequently for work or personal reasons.

Benign paroxysmal positional vertigo (BPPV):

Benign paroxysmal positional vertigo (BPPV) is the most common vestibular disorder. It occurs when tiny calcium particles clump up in the inner ear, causing brief episodes of mild to intense dizziness. Diagnosis is typically through the Dix-Hallpike test, which involves observing involuntary eye movements as the patient's head is moved into specific positions. Treatment often involves maneuvers to move the calcium deposits out of the canal they're affecting.

Vestibular neuritis:

Vestibular neuritis, an inner ear inflammation usually caused by viral infections, leads to sudden, severe vertigo, nausea, and imbalance. Diagnosis typically involves ruling out other causes of these symptoms. Treatments usually focus on alleviating the symptoms and may include medications, vestibular rehabilitation therapy, and lifestyle changes.

 Core: Anatomy and Physiology

Hormones of the Adrenal Glands

JoVE 14987

Adrenal hormones play a pivotal role in maintaining the body's electrolyte balance and orchestrating responses to stress, showcasing the intricate functions of the adrenal cortex and medulla.

The adrenal cortex, a powerhouse of hormone synthesis, generates over two dozen corticosteroid hormones. The zona glomerulosa produces mineralocorticoids, exemplified by aldosterone, influencing the electrolyte composition of body fluids. The synthesis of glucocorticoids such as cortisol and corticosterone occurs in the zona fasciculata. Aptly named for their impact on glucose metabolism, glucocorticoids accelerate glucose synthesis in the liver and promote glycogen formation in muscle and adipose tissue.

Additionally, the zona fasciculata produces small amounts of androgens, contributing to the development of pubic hair in adolescents and supporting muscle mass, blood cell formation, and the sex drive in adult women.

The adrenal medulla, on the other hand, houses chromaffin cells that secrete catecholamines in response to stress. Epinephrine and norepinephrine, continuously stored in vesicles and released at low levels through exocytosis, are crucial in promoting strength and endurance in skeletal muscles. Simultaneously, these hormones facilitate the breakdown of stored fat and glycogen in adipose tissue and the liver, providing the necessary energy for the body's heightened demands during stressful situations. The adrenal hormones finely tune the body's responses to ensure adaptability and survival.

 Core: Anatomy and Physiology

Mesh Analysis

JoVE 15049

Mesh analysis is a valuable method for simplifying circuit analysis using mesh currents as key circuit variables. Unlike nodal analysis, which focuses on determining unknown voltages, mesh analysis applies Kirchhoff's voltage law (KVL) to find unknown currents within a circuit. This method is particularly convenient in reducing the number of simultaneous equations that need to be solved.

A fundamental concept in mesh analysis is the definition of meshes and mesh currents. A mesh is a closed loop within a circuit that does not contain any other loops within it. Each mesh is assigned a mesh current, typically assumed to flow in a clockwise direction within its respective loop.

For mesh analysis to be applicable, the circuit must be planar, meaning it can be drawn on a flat surface without branches crossing one another. Planar circuits are ideal for mesh analysis, as it simplifies the process. The steps involved in mesh analysis are as follows:

  • • Assign mesh currents to each of the "n" independent meshes in the circuit.
  • • Apply KVL to each of the "n" meshes, expressing element voltages in terms of mesh currents using Ohm's law.
  • • Solve the resulting "n" simultaneous equations to obtain the values of the mesh currents.

These mesh currents can then determine various branch currents within the circuit. It is important to note that mesh currents are distinct from branch currents unless a mesh is isolated.

 Core: Electrical Engineering

Voltage Dividers

JoVE 15072

In electrical circuits, resistors can be connected in series, sequentially linked one after the other. In a series configuration, the same current flows through each resistor. Ohm's law is a fundamental principle to understand the behavior of resistors in series. It expresses the voltage across these resistors in terms of the current and resistance.

Kirchhoff's voltage law implies that the sum of the voltages across the resistors in series equals the source voltage. This means that the current in the circuit is equal to the source voltage divided by the total resistance in series. This current expression can be substituted into Ohm's law to determine the voltage across each resistor. This arrangement results in a proportional division of the source voltage among multiple resistors based on their individual resistances. This principle is known as "voltage division," a circuit that exemplifies this principle is called a "voltage divider."

A combination of resistors connected in series can be treated as a single equivalent resistor. The value of this equivalent resistor is simply the sum of the resistances of the individual resistors in series. This concept extends to any number of resistors connected in series.

Equation1

The voltage drop across each resistor in a series circuit is directly proportional to its resistance. As a result, resistors with larger resistance values will experience more significant voltage drops. This principle is at the core of voltage dividers, where each resistor's voltage drop is determined by its resistance relative to the total resistance in the series.

In practice, combining resistors in series or parallel is common, often necessitating a simplified representation of the circuit. This simplification is achieved by combining two resistors at a time. The equivalent resistance of two resistors in a series is the sum of their individual resistances. This technique facilitates the analysis of complex circuits, providing engineers with a valuable tool for designing and troubleshooting electrical systems.

 Core: Electrical Engineering

Inductors

JoVE 15088

An inductor is a passive component built to store energy within its magnetic field. It can be fabricated by coiling a wire around a magnetic core. When current is permitted to flow through this inductor, it is observed that the voltage across the inductor is directly proportional to the time rate of change of the current. Mathematically,

Equation1

This relationship is expressed using the passive sign convention, where 'L' represents the constant of proportionality, also known as the inductance of the inductor.

The inductance of an inductor is influenced by factors such as its size, the materials used, and the construction method. The unit of measurement for inductance is the Henry (H). Inductance is a measure of the ability of a device to store energy in the form of a magnetic field. Also, it is a characteristic that allows an inductor to resist changes in the current flowing through it.

In addition to being referred to as inductors, these components can also be called coils or chokes. The circuit symbol for an inductor is presented in Figure 1.

Figure1

Figure 2 shows a graphical representation of the relationship between voltage and current for an inductor whose inductance does not depend on the current, also known as a linear inductor.

Figure2

A nonlinear inductor, on the other hand, would not have a straight-line plot, as its inductance fluctuates with the current.

Inductors, akin to capacitors, are commercially available in various values and types. In practical applications, inductors typically have inductance values ranging from a few microhenrys, as seen in communication systems, to tens of henrys, as used in power systems. Inductors can either be fixed or variable, and their core can be composed of various materials, including iron, steel, plastic, or even air.

 Core: Electrical Engineering

Phasors

JoVE 15106

Phasors are a powerful mathematical tool used to analyze alternating current (AC) circuits. They provide a complex number representation of sinusoids, with the magnitude of the phasor equating to the amplitude of the sinusoid and the angle of the phasor representing the phase measured from the positive x-axis.

One of the significant benefits of using phasors is that they simplify the analysis of AC circuits by eliminating the time dependence of the current and voltage. This transformation allows an AC circuit to be analyzed as if it were its equivalent direct current (DC) form, making calculations more straightforward.

Phasors can be represented in different forms - rectangular, polar, or exponential - by using Euler's identity, a fundamental formula in complex analysis that establishes a deep relationship between trigonometric and exponential functions.

To obtain the phasor of a sinusoid in sine form, one must first convert it into cosine form and then express it as the real part of a complex number. The phasor of this sinusoid equals the time-independent part of this complex number. Conversely, the sinusoid of a given phasor can be obtained by multiplying the phasor with a time factor and taking its real part.

In a graphical context, phasors can be visualized as rotating vectors, or 'sinors,' spinning in a counterclockwise direction on a complex plane with a constant angular frequency. The diagrams that depict these rotating sinors are known as phasor diagrams.

A key concept in understanding the phasor diagram is that the projection of the rotating sinors onto the real axis represents the sinusoids. This means that the horizontal position of the sinor at any point in time corresponds to the instantaneous value of the sinusoidal function it represents.

 Core: Electrical Engineering

Design Example: Underdamped Parallel RLC Circuit

JoVE 15175

Consider designing an oscillator circuit, a crucial component in various electronic devices and systems. The objective is to create an oscillator circuit with specific characteristics: a damped natural frequency of 4 kHz and a damping factor of 4 radians per second. To accomplish this, a parallel RLC circuit is employed, known for its ability to sustain oscillations at a resonant frequency. In this case, the damping factor is pivotal in achieving the desired performance.

Starting with a fixed resistance of 200 Ω, the necessary values for capacitance and inductance are determined to meet the specifications. First, it is recognized that the damping factor is the reciprocal of twice the product of resistance and capacitance. This insight allows the pinpointing of the required capacitance value for the circuit.

Equation1

Next, attention is turned to the mathematical expression that connects the resonant frequency, damping factor, and damped natural frequency. By manipulating this equation, resonant frequency information can be extracted.

Equation2

The resonant frequency, crucial for oscillator design, is inversely proportional to the square root of the product of inductance and capacitance. Armed with this knowledge and the known values, the necessary inductance to meet the design criteria is computed.

Equation3

These calculations confirm that the damped natural frequency is indeed lower than the resonant frequency, and the inductance value falls below four times the square of the resistance multiplied by capacitance. These conditions affirm that the oscillator circuit has been designed to exhibit the required underdamped oscillations, showcasing proficiency in precise circuit design.

 Core: Electrical Engineering

Impedance Combination

JoVE 15569

Consider a string of christmas lights, each bulb symbolizing an impedance element. In this series configuration, the flow of electric current remains uniform across every component. This behavior aligns with Kirchhoff's Voltage Law (KVL), which asserts that the total impedance in such a setup equals the sum of individual impedances—akin to resistors in series. It follows that the voltage from the power source is distributed proportionally among these components, adhering to the voltage division principle.

Equation1

However, the drawback of this series connection is evident when a single bulb fails, causing an open circuit that interrupts the entire current flow. christmas lights are typically arranged in a parallel configuration to ensure a continuous and steady power supply. This setup guarantees a constant voltage across each bulb, as per Kirchhoff's Current Law (KCL), where the reciprocal of the equivalent impedance equates to the sum of the reciprocals of individual impedances—similar to resistors in parallel. It follows that the equivalent admittance is the sum of the individual admittances.

In this parallel arrangement, the source current divides inversely based on the impedances of the individual bulbs, exemplifying the current division principle. Notably, each bulb establishes an independent pathway to the power source, enabling isolated and uninterrupted current flow.

Equation2

Equation3

Furthermore, in more complex circuits with both series and parallel impedances, the delta-to-wye and wye-to-delta transformations can also be employed for impedance circuits, offering valuable circuit analysis and design tools. These transformations facilitate the conversion between different impedance configurations, enhancing the versatility of impedance-based circuits.

 Core: Electrical Engineering

Zebrafish Breeding and Embryo Handling

JoVE 5150

Zebrafish (Danio rerio) are an important model organism that is particularly valuable for research in developmental biology. Zebrafish are extremely fertile and can produce hundreds of progeny per week, so it is relatively easy to collect a large number of embryos for high sample numbers. Furthermore, zebrafish undergo rapid development and embryos are transparent, allowing for easy visualization of developmental processes.

This video covers the steps required for the collection of newly fertilized zebrafish embryos. A brief overview of zebrafish mating behavior is presented, followed by instructions for setting up crosses in specialized laboratory breeding tanks that allow for controlled mating. Also covered are the conditions required to initiate the release of eggs (known as spawning) the morning after tanks are set. Next, essential techniques for working with embryos are presented, including the inhibition of pigment development with the chemical PTU, and dechorionation: a procedure in which the shell-like membrane surrounding the embryo (the chorion) is removed. Finally, the video concludes with some practical applications of these techniques in developmental research.

 Biology II

Patch Clamp Electrophysiology

JoVE 5202

Neuron cell membranes are populated with ion channels that control the movement of charge into and out of the cell, thereby regulating neuron firing. One extremely useful technique for investigating the biophysical properties of these channels is called patch clamp recording. In this method, neuroscientists place a polished glass micropipette against a cell and apply suction to form a high resistance seal. This process isolates a small “patch” of membrane that contains one or more ion channels. Using an electrode within the micropipette, researchers can “clamp” or control the electrical properties of the membrane, which is important for analysis of channel activity. The electrode also allows for changes in the voltage across the membrane, or the flow of ions through the membrane, to be recorded.

This video begins with an overview of the principles behind patch clamp electrophysiology, an introduction to the necessary equipment, and descriptions of the various patch configurations, including whole cell, cell-attached, perforated, inside-out, and outside-out patches. Next, the key steps of a typical whole-cell patch clamp experiment are outlined, in which a current-voltage (IV) curve is generated. Finally, applications of patch clamp recording are provided to demonstrate how the biophysical properties of ion channels, cell excitability, and neuroactive compounds are evaluated in neurophysiology labs today.

 Neuroscience

An Introduction to Developmental Genetics

JoVE 5325

Development is the complex process through which a single-celled embryo transforms into a multicellular organism. Developmental processes are guided by information encoded in an organism's DNA, and geneticists are trying to understand how this information leads to a fully formed organism.

This video reviews seminal research in the field of developmental biology, including the identification of specific genes that control various embryonic processes. An introduction to the major questions asked by developmental geneticists, and the prominent methods used to answer them, is also provided. Finally, several applications of these prominent methods are discussed, in order to show specific experiments currently being performed in this field.

 Developmental Biology

X-ray Fluorescence (XRF)

JoVE 5498

Source: Laboratory of Dr. Lydia Finney — Argonne National Laboratory

X-ray fluorescence is an induced, emitted radiation that can be used to generate spectroscopic information. X-ray fluorescence microscopy is a non-destructive imaging technique that uses the induced fluorescence emission of metals to identify and quantify their spatial distribution.

 Analytical Chemistry

Expression Profiling with Microarrays

JoVE 5547

Microarrays are important tools for profiling gene expression, and are based on complementary binding between probes that are attached to glass chips and nucleic acids derived from samples. Using these arrays, scientists can simultaneously evaluate the expression of thousands of genes. In addition, the expression profiles of different cells or tissue types can be compared, allowing researchers to deduce how the expression of different genes change during biological processes, and thus gain insight into how the genes may function in pathways or networks.

Here, JoVE explains the principles behind microarrays. This is followed by a general protocol for performing a microarray experiment, and a brief introduction to analyzing microarray data. We end on a discussion of how scientists are currently using microarrays, for example to compare gene expression between different cell types derived from cancerous and non-cancerous tissues, to study important biological problems.

 Genetics

Live Cell Imaging of Mitosis

JoVE 5642

Mitosis is a form of cell division in which a cell’s genetic material is divided equally between two daughter cells. Mitosis can be broken down into six phases, during each of which the cell’s components, such as its chromosomes, show visually distinct characteristics. Advances in fluorescence live cell imaging have allowed scientists to study this process in great detail, providing important insights into the biological control of this process and how it might go wrong in diseases such as cancer.

We begin this video by breaking down the phases of mitosis, and introducing some important considerations for optimal visualization of the process using live cell imaging. We then walk through the steps for running a live cell mitosis imaging experiment and discuss various analysis methods, including the generation of montages, movies, and 3D recreations. Finally, we take a look at how visualizing the mitotic process can be applied to answering questions in cell biology.

 Cell Biology

Dialysis: Diffusion Based Separation

JoVE 5684

Dialysis is a common technique used in biochemistry for separating molecules based on diffusion. In this procedure, a semipermeable membrane allows the movement of certain molecules based on size. This method can be applied to the removal of buffer, known as desalting, or exchanging buffer molecules or ions from a protein solution.

This video covers the principles of dialysis along with a general procedure.  Several applications of dialysis are reviewed, including the removal of gradient reagents following ultracentrifugation, removing detergent after a membrane protein extraction, and the reconstitution of proteins by changing the solution environment.

Biochemical samples typically have high buffer concentrations that can disrupt downstream processing and analysis. Dialysis is a common, inexpensive technique used to separate molecules based on diffusion. The method utilizes a semi-permeable membrane that allows the movement of certain components, based on size. This video will show the concepts of dialysis, a general procedure, and some of its uses in biochemistry.

The most important aspect of dialysis is a semi-permeable membrane, which has pores that impose a molecular weight cut-off, allowing molecules below a certain size to pass through. For example, a 10k membrane will generally retain molecules larger than 10 kilodaltons. However, the molecular weight cutoff is not a discrete or precise boundary. The membrane typically contains a broad range of pore sizes, so a small fraction of molecules near the cutoff may be lost.

Since the molecular weight cutoff is defined using globular proteins, linear molecules of similar mass, like DNA or RNA, may slip through. Membranes are typically chosen one half to one third the molecular weight of the desired molecule.

To perform the procedure, a sample is placed into the membrane, which is in turn added to a large volume of solution, called the dialysate. Over time, smaller molecules will diffuse freely across the membrane between the sample and the dialysate, while the larger biomolecules are held within. Dialysis is a slow process. It is common to allow it to run overnight, or even across multiple days.

If the dialysate is pure water, the overall buffer concentration will decrease, a process known as desalting. If the solution contains other small particles, some will move into the sample, leading to buffer exchange. Because dialysis is an equilibrium process, the dialysate can be refreshed multiple times to further displace small molecules. Once the process is complete the sample is re-collected for further processing.

Now that you've seen the basics of dialysis, let's take a look at a general procedure.

Before beginning the procedure, the membrane is presoaked in dialysate. This makes it easier to use, and removes any preservatives. Once ready, the sample is collected, typically with a syringe and is then added to the dialysis container. This can be bare tubing, or contained within a cassette. Excess air is removed from the dialysis setup to maximize the sample's surface area with the membrane. The setup is then placed into the dialysate with stirring to maximize the diffusion. It should float to not inhibit stirring.

The dialysate is changed at relevant intervals as equilibrium between sample and dialysate is reached. After the last change, the reaction is typically left to run overnight. After a sufficient time period, the buffer-free or -exchanged sample is removed from the cassette. Once collected, the sample can be analyzed or further processed, depending on the nature of the experiment.

Now that we've looked at a general dialysis procedure, let's see some of the ways this technique is used in biochemistry.

Density gradients are a common way to separate complex biological samples. This concept relies on the distribution on small particles, typically sucrose or cesium chloride ions. Once complete, these reagents typically need to be removed before the collected sample can be processed. Dialysis makes it possible to utilize the purified sample for future analysis.

Certain proteins are found within a cell's lipid bilayer, and are usually studied by interspersing them into spherical lipid vesicles known as liposomes. The proteins and lipids are first extracted with a detergent. Dialysis can be used to slowly remove the detergent, forming proteoliposomes.

After purification, some proteins are misfolded, or denatured, leading to a loss in functionality. The compounds that cause these changes in structure can be removed with dialysis, leading to the reformation of functioning analytes.

You've just watched JoVE's video on dialysis. You should now understand this diffusion-based method, a simple experimental procedure, and the use of this technique.

Thanks for watching!

 Biochemistry

Raman Spectroscopy for Chemical Analysis

JoVE 5701

Source: Laboratory of Dr. Ryoichi Ishihara — Delft University of Technology

Raman spectroscopy is a technique for analyzing vibrational and other low frequency modes in a system. In chemistry it is used to identify molecules by their Raman fingerprint. In solid-state physics it is used to characterize materials, and more specifically to investigate their crystal structure or crystallinity. Compared to other techniques for investigating the crystal structure (e.g. transmission electron microscope and x-ray diffraction) Raman micro-spectroscopy is non-destructive, generally requires no sample preparation, and can be performed on small sample volumes.

For performing Raman spectroscopy a monochromatic laser is shone on a sample. If required the sample can be coated by a transparent layer which is not Raman active (e.g., SiO2) or placed in DI water. The electromagnetic radiation (typically in the near infrared, visible, or near ultraviolet range) emitted from the sample is collected, the laser wavelength is filtered out (e.g., by a notch or bandpass filter), and the resulting light is sent through a monochromator (e.g., a grating) to a CCD detector. Using this, the inelastic scattered light, originating from Raman scattering, can be captured and used to construct the Raman spectrum of the sample.

In the case of Raman micro-spectroscopy the light passes through a microscope before reaching the sample, allowing it to be focused on an area as small as 1 µm2. This allows accurate mapping of a sample, or confocal microscopy in order to investigate stacks of layers. Care has to be taken, however, that the small and intense laser spot does not damage the sample.

In this video we will briefly explain the procedure for obtaining a Raman spectra, and an example of a Raman spectrum captured from carbon nanotubes will be given.

 Analytical Chemistry

Chemical Bonds

JoVE 11679

Atoms participate in a chemical bond formation to acquire a completed valence-shell electron configuration similar to that of the noble gas nearest to it in atomic number. Ionic, covalent, and metallic bonds are some of the important types of chemical bonds. Bond energy and bond length determine the strength of a chemical bond.

Types of Chemical Bonds

An ionic bond is formed due to electrostatic attraction between cations and anions. Often, the ions are formed by the transfer of electrons from one participating atom to the other. However, these bonds do not have a defined directionality because the electrostatic force of attraction is distributed uniformly throughout the three-dimensional space.

A covalent bond is a chemical bond formed by the sharing of electron pairs between adjacent atoms. The shared pair of electrons is called the bonding pair. Covalent bonds are directional in nature.

A metallic bond is formed between two metal atoms. Metallic bonding is described by the “Electron Sea model”. Based on the low ionization energies of metals, the model states that metal atoms lose their valence electrons easily and become cations. These valence electrons create a pool of the delocalized electrons surrounding the cations over the entire metal.

Bond Energies and Bond Length

The strength of a covalent bond is measured by the energy required to break it—that is, the energy necessary to separate the bonded atoms. Separating any pair of bonded atoms requires energy. The stronger a bond, the greater the energy required to break it.

The energy required to break a specific covalent bond in one mole of gaseous molecules is called the bond energy or the bond dissociation energy. The bond energy for a diatomic molecule is defined as the standard enthalpy change for the endothermic reaction. Molecules with three or more atoms have two or more bonds. The sum of all bond energies in such a molecule is equal to the standard enthalpy change for the endothermic reaction that breaks all the bonds in the molecule.

The strength of a bond between two atoms increases as the number of electron pairs in the bond increases. Generally, the greater the number of bonds between two atoms, the shorter the bond length and the greater the bond strength. Thus, triple bonds are stronger and shorter than double bonds between the same two atoms; likewise, double bonds are stronger and shorter than single bonds between the same two atoms. When one atom bonds to various atoms in a group, the bond strength typically decreases as we move down the group.

This text is adapted from Openstax, Chemistry 2e, Section 7.1: Ionic BondingOpenstax, Section 7.2: Covalent BondingSection 10.5: The Solid State of Matter, and Section 7.5. Bond Strength: Covalent Bonds.

 Core: Organic Chemistry

Physical Properties of Ethers

JoVE 11722

Overview

An ether molecule has a net dipole moment due to the polarity of C–O bonds. Subsequently, boiling points of ethers are lower than those of alcohols of comparable molecular weight and slightly higher than those of hydrocarbons of comparable molecular weight (Table 1).

Ethers can act as hydrogen bond acceptors, making them more water-soluble than hydrocarbons, but since ethers cannot act as hydrogen bond donors, they are much less soluble in water than alcohols. Ethers are considered good solvents because of their ability to form hydrogen bonds with other molecules, combined with the London Dispersion forces between the alkyl groups bonded to oxygen. Ethers have high volatility and can quickly evaporate during the isolation of reaction products.

Table 1. Comparison of Boiling Points of Ethers, Alcohols, and Hydrocarbons

Name  Structural Formula Molecular weight (g/mol)  bp (°C)
Dimethyl ether Figure1 46 −25
Ethanol Figure1 46 78
Propane Figure1 44 −42
Diethyl ether Figure1 74 35
1-Butanol Figure1 74 118
Pentane Figure1 72 36

 Core: Organic Chemistry

Relative Stabilities of Alkenes

JoVE 11768

The relative stability of alkenes can be determined by comparing their heats of hydrogenation. The lower heat of hydrogenation indicates the more stable alkene.  The three main factors determining the relative stability of alkenes are i) the number of substituents attached to the double-bond carbon atoms, ii) hyperconjugation, and iii) the stereochemistry of the double bond.

  1. Number of substituents across the double bond: An alkene with two smaller substituents is more stable than its isomer having one large substituent. For example, 2-butene is more stable compared to 1-butene. The highly substituted alkenes have a higher ratio of sp2sp3 bonds, which are lower in energy and are stronger as compared to the sp3sp3 bonds. Thus, a tetrasubstituted alkene is more stable than a tri-, di-, or monosubstituted alkene.
  2. Hyperconjugation: Hyperconjugation is a stabilizing interaction of the delocalized electron density between the carbon–carbon π bond and the adjacent carbon–hydrogen σ bonds on the substituent. Thus, a higher number of alkyl substituents across the double bond suggests greater hyperconjugation, resulting in a more stable alkene.
  3. Stereochemistry: The spatial arrangement of the substituents also contributes to the stability of alkenes. The cis isomer exhibits steric strain because of the crowding of the substituents on the same side of the double bond and therefore is less stable compared to the trans isomer.

 Core: Organic Chemistry

Crown Ethers

JoVE 11787

Crown ethers are cyclic polyethers that contain multiple oxygen atoms, usually arranged in a regular pattern. The first crown ether was synthesized by Charles Pederson while working at DuPont in 1967. For this work, Pedersen was co-awarded the 1987 Nobel Prize in Chemistry. Crown ethers are named using the formula x-crown-y, where x is the total number of atoms in the ring and y is the number of ether oxygen atoms. The term 'crown' refers to the crown-like shape that these ether molecules take. A significant feature of crown ethers is that they form complexes with specific alkali metal cations. The oxygen atoms of crown ethers together form an internal cavity into which the electron lone pairs effectively coordinate the metal ions. The choice of the metal ion depends on the diameter of the ether's internal cavity compared to the diameter of the metal ion. Consequently, crown ethers serve as effective solvating agents for solubilizing inorganic salts in organic solvents. For example, KF would not dissolve in benzene by itself, but the use of 18-crown-6 generates a complex with potassium ion, which dissolves in benzene.

Figure1

The result is a solution containing unsolvated fluoride ions, free to participate in nucleophilic substitution reactions. Typically, the strong interaction between fluoride ions and polar solvents makes it challenging to free up fluoride anions in a nonpolar solvent. However, crown ether increases the nucleophilic strength of the fluoride anion by making it available to participate in an SN2 reaction. Overall, the role of crown ether is to sequester the cation, leaving the anion to function as a better nucleophile.

 Core: Organic Chemistry

Protection of Alcohols

JoVE 11927

This lesson delves into the concept of protection and deprotection of a functional group fundamental to synthetic organic chemistry. These phenomena are explained in the context of aliphatic and aromatic alcohols.

Protection

It defines a protecting group as the masking agent to make the more reactive species inert to a given set of conditions. This concept is depicted via the illustration of liquid flow through different outlets in an assembly of pipes. The analogy helps to understand the role of a protecting group in reaction selectivity, as in the case of the organolithium alkylation of a halide in the presence of a competing acidic alcohol group. The example shows how protection of the alcohol group helps to achieve the alkylation of the halide. Popular protecting groups for alcohols include the trialkylsilyl family for nucleophiles or carbon and nitrogen bases and the tetrahydropyranyl (THP) group for strong bases. In the former example, the halide of the trialkylsilyl derivative reacts with the alcohol in the presence of a nucleophilic catalyst to generate a trialkylsilyl ether.

Deprotection

Every protection is followed by deprotection after the intended reaction. The deprotection restores the system to its native state. In protection with trialkylsilyl groups, deprotection is achieved using fluoride salts like tetra-n-butylammonium fluoride (TBAF) that are soluble in organic solvents. Here, the re-protonation of the oxygen regenerates the native alcohol. In the case of protection with THP, deprotection is achieved using acid hydrolysis.

Principle of design

The lesson also elucidates the principles behind the design of a protecting group using an illustration of a house under varying external weather conditions. It demonstrates the selectivity offered by a protecting group in a specific environment. For instance, THP protects alcohol from strong bases. The acetal formed in this case is stable towards bases but susceptible to acid hydrolysis.

Apart from the reaction conditions, the reactivity of the molecule to be protected also plays a key role in designing a suitable protecting group. For example, methyl ethers’ ability to protect phenols is found inappropriate for aliphatic alcohols. Here, the stability of the corresponding leaving groups during deprotection plays a key role. For example, the alkoxides, unlike phenoxides, are poor leaving groups for deprotection with hydrogen bromide.

The following table summarizes the various protection/deprotection groups for different types of alcohols and related conditions:

Protecting group Structure Protects From Protection Deprotection
Trialkylsilyl (R3Si–),
e.g., TBDMS
Me3Si–OR
(Me3C)Me2Si–OR
Alcohols
(OH in general)
Nucleophiles,
C or N bases
R3SiCl,
base
H+, H2O,
or F
Tetrahydropyranyl
(THP)
Figure1 Alcohols
(OH in general)
Strong bases 3,4-Dihydropyran,
H+
H+, H2O
Benzyl ether
(OBn)
Figure2 Alcohols
(OH in general)
Almost everything NaH, BnBr H2, Pd/C,
or HBr
Methyl ether
(ArOMe)
Figure3 Phenols
(ArOH)
Bases NaH, MeI, or
(MeO)2SO2
BBr3, HBr, HI,
Me3SiI

 Core: Organic Chemistry

Electrical Synapses

JoVE 12181

Electrical synapses found in all nervous systems play important and unique roles. In these synapses, the presynaptic and postsynaptic membranes are very close together (3.5 nm) and are actually physically connected by channel proteins forming gap junctions.

Gap junctions allow the current to pass directly from one cell to the next. In contrast, in the chemical synapse, the neurotransmitters carry the information through the synaptic cleft from one neuron to the next. They consist of two hexameric connexin hemichannels or connexons contributed by each of the adjacent cells. These hemichannels make contact between the two cell membranes by forming a continuous bridge between the cytoplasm of the two connecting cells. The opening of the connexon pore is more like the shutter of a camera where the connexins in the hemichannel rotate slightly with respect to one another for ion passage. In addition to the ions, other molecules, such as ATP, can also diffuse through the large gap junction pores.

Signaling in electrical synapses is virtually instantaneous. Some electrical synapses are bidirectional too. Electrical synapses are not easily blocked and are important for synchronizing the electrical activity of a group of neurons. For example, electrical synapses in the thalamus are thought to regulate slow-wave sleep, and disruption of these synapses can cause seizures. In the intestinal smooth muscle cells, electrical synapses provide electrical rhythmicity contributing to peristaltic intestinal activity vital for normal functioning of the gastrointestinal tract.

 Core: Cell Biology

Basicity of Heterocyclic Aromatic Amines

JoVE 12527

Heterocyclic amines, where the N atom is a part of an alicyclic system, are similar in basicity to alkylamines. Interestingly, the heterocyclic amine having a nitrogen atom as part of an aromatic ring has much less basicity than its corresponding alicyclic counterpart. For this reason, as presented in Figure 1, piperidine (pKb = 2.8) is significantly more basic than pyridine (pKb = 8.8).

Figure1

Figure 1. The comparison of the basicity of piperidine and pyridine.

This difference in basicity may be attributed to the state of hybridization of orbitals containing a lone pair of electrons on the N atom, as depicted in Figure 2. In the case of piperidine, the lone pair resides in an sp3-hybridized orbital having lower s character, making the lone pair more available for exhibiting basicity towards acid. On the other hand, in piperidine, the lone pair resides in an sp2-hybridized orbital containing a much higher s character. Consequently, the lone pair is more tightly bound to the aryl ring and less available for exhibiting basicity towards acid.

Figure2

Figure 2. The effect of the hybridized orbitals on the basicity.

Pyrrole is much less basic than pyridine, having a pKb value of 15. The lone pair of electrons on the N atom shown in Figure 3 resides in a pure p orbital and are perfectly aligned with the p orbitals of C atoms, which participate in the ring’s aromaticity. Therefore, the lone pair of electrons on the N atom of pyrrole is delocalized by resonance throughout the ring. Conversely, the lone pair on the N atom of pyridine being in an sp2 hybridized orbital is aligned perpendicular to the other sp2 orbitals of C atoms taking part in resonance. So, the lone pair of electrons on the N atom of pyridine is much more available than that in pyrrole, resulting in higher basicity of pyridine.

Figure3

Figure 3. The Lewis structure of pyrrole.

Imidazole contains two N atoms in a five-membered ring and is a vital heterocycle found in many proteins. This heterocyclic amine with a pKb of 7 is more basic than pyridine by a factor of 100. The N atom resembling that in pyrrole is non-basic. The other N atom is basic and abstracts H from acid giving rise to its conjugated acid, which is stabilized by the resonance, as illustrated in Figure 4. This resonance stabilization of the conjugate base results in the increased basicity of imidazole relative to pyridine.

Figure4

Figure 4. The resonance stabilization in imidazolium ions.

Basicity of amines provides a valuable practical tool to separate amines from a mixture containing other neutral compounds. This is achieved by dissolving the impure mixture in ether and shaking it with water in a separatory funnel. After the two layers separate, draining the water layer removes the salt and most inorganic impurities. When dilute aqueous acid is added to the organic layer, amines are selectively protonated into the corresponding acid and dissolve in the aqueous layer. Draining the organic layer removes the neutral organic impurity. Slow basification of the aqueous layer again regenerates the free amine, which could be extracted using a fresh volume of ether. The subsequent evaporation of ether yields a pure amine.

 Core: Organic Chemistry

Role of Septins

JoVE 12600

Septins are the recently discovered fourth major protein component of the cytoskeleton, along with microfilaments, microtubules, and intermediate filaments. These proteins can associate with other cytoskeletal filaments and carry out varied roles or can be free-floating in the cytoplasm.

Cellular Functions of Septins

Recent studies have revealed the multifaceted roles of septins in various cellular processes such as cytokinesis, ciliogenesis, and neurogenesis. Septins act as scaffolds and promote protein-protein interactions; for example, in S. cerevisiae, septin scaffolding protein helps recruit different filament proteins to form contractile rings during cytokinesis.

Septin was first found in the sperm flagella's annulus, which compartmentalizes the anterior and posterior tail regions. Mutation in the septin-forming genes like SEPT12 affects sperm motility and its structural integrity. Other septins like SEPT2 are involved in forming a diffusion barrier between the cilia and the cytoplasm.

Septins also play a role in stabilizing the membrane by binding and modifying the membrane and action-myosin interactions.

Diseases Related to Septins

Studies have shown septins are associated with proteins involved in neurodegenerative diseases like Parkinson's and Alzheimer's. In Parkinson's disease, septin interacts with parkin, an E3 ubiquitin ligase. Septins' mutation has also been linked with tumorigenesis or cancer. Studies on cancerous cells suggest septin mutations are also related to resistance to cells, metastasis, proliferation, and angiogenesis.

 Core: Cell Biology

GPCRs Regulate Adenylyl Cyclase Activity

JoVE 13324

Some GPCRs transmit signals through adenylyl cyclase (AC), a transmembrane enzyme. AC helps synthesize second messenger cyclic adenosine monophosphate (cAMP). AC catalyzes cyclization reaction and converts ATP to cAMP by releasing a pyrophosphate. The pyrophosphate is further hydrolyzed to phosphate by the enzyme pyrophosphatase, which drives cAMP synthesis to completion. However, cAMP is rapidly degraded to 5′ AMP by the enzymes phosphodiesterase (PDE), preventing overstimulation of cells.

Two types of heterotrimeric G proteins regulate AC.

  1. The stimulatory G protein (Gs) binds and activates the AC, increasing cAMP levels.
  2. The inhibitory G protein (Gi) binds and inactivates AC, lowering cAMP levels.

AC consists of a small N-terminal domain, two repeated transmembrane, and a cytoplasmic domain. The cytoplasmic domain forms the catalytic core and consists of the G protein binding site. The binding of Gs protein to AC help opens its catalytic core, facilitating the binding of substrate, ATP. In contrast, binding of Gi protein leads to reorientation of catalytic residues, thereby closing the catalytic core, which prevents ATP binding.

Many pathogenic bacterial toxins disrupt this regulation mechanism of adenylyl cyclases in the host cells. For example, Vibrio cholerae enter intestinal epithelial cells and release cholera toxins. The toxin induces ADP ribosylation of the G⍺s  subunit of the stimulatory G protein.  The addition of ADP-ribose inhibits intrinsic GTPase activity of the G⍺s  subunit. Thus, the modified G⍺s  subunit remains constitutively active, activating adenylyl cyclase without external stimulation. Continued AC activation prolongs cAMP synthesis. The resulting increase in cAMP concentration causes water and electrolyte efflux from the intestinal cells to the lumen. This leads to severe diarrhea and excessive dehydration, characteristic of cholera.

Another pathogen infecting the airways, Bordetella pertussis, produces a toxin that catalyzes ADP ribosylation of the G⍺i subunit of the inhibitory G protein. This modification prevents GDP from leaving the G⍺i  subunit, thereby keeping Gi protein in the inactive state. Gi protein cannot bind and inhibit adenylyl cyclase activity. Thus, cAMP released in the airway epithelial cells leads to loss of fluids and electrolytes along with enhanced mucus secretion, causing whooping cough.

 Core: Cell Biology

Role of Microtubules in Cell Wall Deposition

JoVE 13368

Microtubules are small hollow tubes in eukaryotic cells. The cell wall microtubules are polymerized dimers of two globular proteins, α-tubulin and β-tubulin, two globular proteins. With a diameter of about 25 nm, microtubules are the widest components of the cytoskeleton. They help the cell resist compression and provide a track along which vesicles move through the cell or pull replicated chromosomes to opposite ends of a dividing cell. Microtubules go through quick cycles of disassembly and reassembly.

In plant cells, cortical microtubules are known to play a role in depositing cellulose in the cell wall. They influence cellulose deposition by either moving microtubules or directly changing the orientation of newly synthesized cellulose microfibrils. Microfibrils stabilize the boundaries of specialized plasma membrane domains that force nascent cellulose chains into a parallel alignment through glucan chain polymerization and chain crystallization. The glucose monomers form hydrogen bonds, thus holding the cellulose chains firmly together, forming oriented microfibrils. This imparts rigidity in the cell wall during cellulose deposition.

Adapted from section 4.5 cytoskeleton, openstax AP biology, section 6.4 Prokaryotic cell division, Openstax concepts of biology, section 3.2 Carbohydrates, Openstax biology 2 e.

 Core: Cell Biology

Peptide Identification Using Tandem Mass Spectrometry

JoVE 13385

Tandem mass spectrometry, also known as MS/MS or MS2, is an analytical technique that employs two mass analyzers. Essentially it is a series of mass spectrometers that helps isolate a particular biomolecule and then helps study its chemical properties.

This technique helps gather information regarding the protein from which the peptide was obtained and to study the peptides’ amino acid sequence. Identifying peptides from a complex mixture is an important component of the growing field of proteomics.

The different peptides obtained after enzymatic digestion can be further separated as much as possible as per the physical size and/or chemical properties using sophisticated instruments like gel electrophoresis or liquid chromatography.

This first stage of MS/MS allows the peptides to separate based on their mass-to-charge ratio, followed by breaking or fragmenting the selected peptide ion in the collision cell. The second mass analyzer helps build the fragmentation pattern to determine the sequence or identify the protein.

Studying the results and deducing the peptide sequence is extremely important once the spectrum is obtained. For this, numerous protein database search algorithms and bioinformatics tools help in sample elucidation from the obtained spectrum. In the case of an unknown protein, the obtained spectrum shows numerous overlapping fragments. However, as the spectral pattern is unique for a given protein, the analysis software compares the obtained spectrum with a database of known peptide sequences, thus elucidating the unknown protein from the overlapping fragments.

Different analyzer combinations can be used to create hybrid MS/MS instruments and thus increase the sensitivity of the results. For example, quadrupole time-of-flight (QTOF) is a combination of quadrupole and time-of-flight mass analyzers. The triple quadrupole mass spectrometer contains two quadrupole mass analyzers separated by a non-mass resolving quadrupole collision cell for fragmentation.

In newborn screening, the tandem MS uses dried blood-spot samples to help in a comprehensive assessment of inborn metabolic disorders like phenylketonuria, sickle cell disease, etc. The early identification helps substantial improvement in health outcomes.

 Core: Cell Biology

Transmission Electron Microscopy

JoVE 13401

In 1931, physicist Ernst Ruska—building on the idea that magnetic fields can direct an electron beam just as lenses can direct a beam of light in an optical microscope—developed the first prototype of the electron microscope. This development led to the development of the field of electron microscopy. In the transmission electron microscope (TEM), electrons are produced by a hot tungsten element and accelerated by a potential difference in an electron gun, which gives them up to 400 keV in kinetic energy. After leaving the electron gun, the electron beam is focused by electromagnetic lenses (a system of condensing lenses) and transmitted through a specimen sample to be viewed. The image of the sample is reconstructed from the transmitted electron beam. The magnified image may be viewed either directly on a fluorescent screen or indirectly by sending it, for example, to a digital camera or a computer monitor.

The entire setup consisting of the electron gun, the lenses, the specimen, and the fluorescent screen are enclosed in a vacuum chamber to prevent energy loss from the beam. Modern high-resolution models of a TEM can have resolving power greater than 0.5 Å and magnifications higher than 50 million times. For comparison, the best resolving power obtained with light microscopy is currently about 97 nm.

A limitation of the TEM is that the samples must be about 100 nm thick, and biological samples require a special preparation involving chemical “fixing” to stabilize them for ultrathin slicing. To overcome the limitations, several advancements in TEM techniques, such as cryo-TEM, have been made to get rid of artifacts and allow for direct sample imaging without the sample-damaging stages of sample fixation and dehydration.  

This text is adapted from Openstax, University Physics Volume 3, Chapter 6 Photons and Matter Waves, Section 6.6: Wave-Particle Duality.

 Core: Cell Biology

Mitosis and Cytokinesis

JoVE 13417

In eukaryotes, the cell division cycle is divided into distinct, coordinated cellular processes that include cell growth, DNA replication/chromosome duplication, chromosome distribution to daughter cells, and finally, cell division. The cell cycle is tightly regulated by its regulatory systems as well as extracellular signals that affect cell proliferation.

The processes of the cell cycle occur over approximately 24 hours (in typical human cells) and in two major distinguishable stages. The first stage is DNA replication, during the S phase of interphase. The second stage is the mitotic (M) phase, which involves the separation of the duplicated chromosomes into two new nuclei (mitosis) and cytoplasmic division (cytokinesis). The two phases are separated by intervals (G1 and G2 gaps), during which the cell prepares for replication and division.

The Process of Mitosis

Mitosis can be divided into five distinct stages—prophase, prometaphase, metaphase, anaphase, and telophase. Cytokinesis, which begins during anaphase or telophase (depending on the cell), is part of the M phase but not part of mitosis.

Prophase

As the cell enters mitosis, its replicated chromosomes begin to condense and become visible as threadlike structures with the aid of proteins known as condensins. The mitotic spindle apparatus begins to form between the centrosomes—which were duplicated during the S phase—and migrate to opposite poles of the cell. The spindle is made up of filamentous structures called microtubules that are comprised of tubulin protein monomers. Spindle microtubules start extending towards the condensed chromosomes. The nucleolus, a component of the nucleus that produces ribosomes, vanishes, indicating the impending breakdown of the nucleus.

Prometaphase

During prometaphase, the microtubule filaments from the spindle apparatus continue to grow, and the chromosomes finish condensing. The nuclear envelope completely breaks down, releasing the chromosomes. Some of the microtubules attach to the released chromosomes, binding at a protein structure called the kinetochore that is present on the centromere of each pair of sister chromatids. Spindle microtubules from opposite poles attach to the kinetochores and capture the condensed sister chromatid pairs. Spindle microtubules that do not attach to chromosomes—polar and astral microtubules—help push the spindles apart and anchor the spindle poles to the cell membrane.

Metaphase

The spindle microtubules align each pair of the fully condensed sister chromatids along the equator of the cell—at the metaphase plate. The cell is now ready to divide.

Anaphase

The microtubules from opposite spindle poles, which are attached to the kinetochore structure, shorten and separate the sister chromatids at the centromere. The cohesion proteins that hold the chromatids together now break down. The shortening kinetochore microtubules cause each chromatid of the pair—now called chromosomes—to migrate to an opposite pole.

Telophase

Once the chromosomes reach opposite poles of the cell, they decondense and uncoil to form chromatin. The spindle microtubule filaments depolymerize into their tubulin monomers, which are then utilized as cytoskeletal elements in daughter cells. Nuclear envelopes reassemble around each set of chromosomes.

Cytokinesis

During cytokinesis in animal cells, actin filaments form a contractile ring in the plasma membrane to create a cleavage furrow, which eventually pinches the cell into two. In plant cells, vesicles from the Golgi apparatus carrying glucose, enzymes, and structural proteins join to form a new cell plate at the location of the former metaphase plate. The growing cell plate fuses with the plasma membranes on each side, eventually forming a new cell wall that divides the cell into two.

Mitosis is now complete, generating two daughter cells that are identical to the parent cell. In most human cells, mitosis accounts for about one hour of the approximately 24-hour cell cycle.

 Core: Cell Biology

Role of Ephrin-Eph Signalling in Intestinal Stem Cell Renewal

JoVE 13465

Erythropoietin-producing hepatocellular carcinoma receptor (Eph) and its ligand, Eph receptor-interacting protein (Ephrin) were first discovered in the human carcinoma cell line, hence the name. Ephrin-Eph interaction guides cells to reach their appropriate location in adult tissues. They also play an essential role in the immune system by helping in immune cell migration, adhesion, and activation. Based on their structure and function, Eph is divided into two classes — EphA and EphB. Similarly, Ephrin is also classified into two types — EphrinA and EphrinB. EphrinA ligands are attached to the plasma membrane through glycosylphosphatidylinositol (GPI), and EphrinB has a transmembrane domain with a cytoplasmic tail. When Ephrin binds to Eph, bidirectional signaling is activated. A forward signal in Eph cells and a reverse signal in Ephrin cells.

As the intestinal cells continuously renew themselves, the Ephrin-Eph signal plays an important role in properly segregating cells and thereby maintaining the intestinal architecture. The Wnt signal triggers EphB2 expression in intestinal stem cells (ISCs) and EphB3 expression in Paneth cells. Mutations in EphB3 can cause misplacement of Paneth cells. In one study with EphB3 mutant mice, the Paneth cells migrated to the villi instead of the crypts’ base.

EphB is a tumor suppressor, and loss of this receptor leads to intestinal tumors. During wound repair, Eph receptor and Ephrin ligands are upregulated in the wounded cells. The resultant cell signaling causes cleavage of adherens junctions between the cells, allowing them to migrate to the wound and help in the healing process.

 Core: Cell Biology

Chromatin Modification in iPS Cells

JoVE 13482

Chromatin modification alters gene expression; therefore, scientists can add histone-modifying enzymes, histone variants, and chromatin remodeling complexes to somatic cells to aid reprogramming into pluripotent stem (iPS) cells.

Compact chromatin makes reprogramming difficult. Enzymes, such as histone demethylases and acetyltransferases, are often added during reprogramming to loosen the chromatin, making the DNA more accessible to transcription factors. Molecules that inhibit histone deacetylases or histone methyltransferases are added to increase the reprogramming efficiency. Similar to histone methylation, DNA methylation also causes chromatin compaction. Inhibitors of DNA methylases help loosen the chromatin and allow the expression of genes essential for pluripotency.

Histone variants can also be added to alter the gene expression pattern. Variants, such as H2AZ and H3.3, change gene expression because they have different amounts of DNA wound around them, allowing specific genes to be more accessible. Additionally, H2AZ often has increased acetylation, enabling more transcription factors to bind to DNA and enhancing reprogramming.

 Core: Cell Biology

Identifying Statistically Significant Differences: The F-Test

JoVE 14515

The F-test is used to compare two sample variances to each other or compare the sample variance to the population variance. It is used to decide whether an indeterminate error can explain the difference in their values. The underlying assumptions that allow the use of the F-test include the data set or sets are normally distributed, and the data sets are independent of each other. The test statistic F is calculated by dividing one variance by another. In other words, the square of one standard deviation by the square of the other. To obtain a value of one or greater than one as the result of the quotient, the larger value is always divided by the smaller value.

The null hypothesis of the F-test states that the ratio is equal to 1. After calculating the test statistic, it is compared to the tabulated critical F values at a chosen confidence level and the appropriate degree of freedom. The null hypothesis is rejected if the test statistic F is smaller than the tabulated F value. In that case, the difference from the desired value of unity–if any–is justified by an indeterminate error, and we state that the variations are not significantly different.

 Core: Analytical Chemistry

Ladder Diagrams: Acid–Base Equilibria

JoVE 14531

Understanding the chemistry between the reagents is necessary for performing any experiment. To this end, scientists have designed a tool called a ladder diagram, which is a graphical representation that helps illustrate the chemistry of a system.

A ladder diagram for acid-base equilibria consists of a vertical axis that represents pH and horizontal bars (steps on the ladder) that help position all the pKa values in the system. At equilibrium, the pH value of the system corresponds to one of the pKa values, which divide the system into more acidic and more basic regions. At pH values higher or lower than any given pKa value, i.e. when the system is not at equilibrium, the dominant species will correspond to the one written in that region of the diagram.

For example, the ladder diagram of the HF and F acid-base equilibria system shows a horizontal line at pH 3.17, which is the pKa value of HF. At pH values above 3.17,  F predominates, whereas at pH values below 3.17, the HF concentration is higher.

The ladder diagram of this system can also be used to understand the effect of pH on the solubility of CaF2. The solubility of CaF2 can be increased by converting F into HF. In contrast, its solubility decreases if F dominantes. From the ladder diagram, it can be understood that pH values above 3.17 allow F to dominate, thereby decreasing the solubility of CaF2

 Core: Analytical Chemistry

Buffers: Buffer Capacity

JoVE 14547

Buffer capacity is the quantitative measure of a buffer to resist the change in pH. As shown in the following equation, the buffer capacity, denoted by 'beta', is expressed as the number of moles of acid or base needed to change the pH of a one-liter buffer solution by 1 unit. Here, Ca and Cb indicate the number of moles of acid and base, respectively. Note that dpH represents the change in pH.

In the graph, pH is plotted as a function of the number of moles of base (Cb) added to a weak acid with pKa equal to 5. The curve's derivative yields another plot that depicts the buffer capacity versus pH. This plot demonstrates that the buffer capacity is highest when pH equals pKa, where the solution contains an equal concentration of the weak acid and its conjugate base. For this weak acid, the buffer capacity is highest at pH 5, where the resistance to pH changes is the highest. Typically, the buffer of choice should have a pKa value that is within plus or minus 1 unit of the desired pH. In addition to the pKa value, buffer capacity also depends on the concentration of the weak acid and its conjugate base in the solution. The higher the buffer species' concentration, the higher the buffer capacity.

 Core: Analytical Chemistry

Colloidal precipitates

JoVE 14586

The high insolubility of some precipitates can result in an unfavorable relative supersaturation. This can lead to colloidal particles with a large surface-to-mass ratio, where adsorption is promoted. For instance, in the precipitation of silver chloride, silver ions are adsorbed on the surface of the colloidal particles, forming a primary layer. This layer attracts ions of opposite charge (such as nitrate ions), forming a diffuse secondary layer of adsorbed ions. This electric double layer prevents colloidal particles from colliding and coagulating into larger particles and stabilizes the suspension.

The coagulation of particles in a colloidal suspension can be enhanced by heating with stirring. This decreases adsorption and increases the kinetic energy of particles to overcome the electrostatic repulsion, enabling coagulation. Alternatively, the addition of an electrolyte can shrink the electrical double layer. At a critical coagulation concentration of electrolyte, the particles can coalesce spontaneously.

Following filtration of a coagulated colloid, washing with pure solvent can decrease the electrolyte concentration below the coagulation value, causing the particles to revert to their dispersed state. This process is called peptization and can be prevented by washing with a non-interfering electrolyte that can be removed by volatilization, such as nitric acid for silver chloride.

 Core: Analytical Chemistry

Coagulation

JoVE 14619

Colloidal solids are solid particles suspended in solution. They are usually negatively charged, attracting a compact primary layer of positively charged ions, which attract more counterions to form an electrical double layer. Electrostatic repulsion between the charged double layers prevents the particles from colliding, stabilizing the colloids. These solids are often undesirable because they can contain toxins that are difficult to remove. Coagulation is a technique that helps aggregate and remove colloids, and it is applicable to wastewater purification.

Metal coagulants are one type of coagulants that can be used in this purification process. When metal coagulants, such as aluminum or ferric salts, are added to the colloidal solution, they react with the bicarbonates in wastewater to yield the respective hydroxides. These hydroxides provide highly charged counterions that can be adsorbed on the suspended colloids, neutralizing their charge. Rapid mixing ensures the dispersal of these highly charged ions and therefore the frequency of formation of neutralized colloids. At the same time, mixing promotes collisions between the neutral colloids, leading to their coalescence as a result of the hydrophobic effect. Upon continued mixing, the particles come together and form slightly larger but still microscopic particles called microflocs. When microflocs reach optimal size and weight and become macroflocs, they can be removed by sedimentation or filtration to obtain pure water.

 Core: Analytical Chemistry

Muscle Stimulation Frequency

JoVE 14848

The contraction strength of muscles is regulated by motor neurons, which modulate the frequency of action potentials dispatched to the motor units based on the body's requirements. This process of varying the muscle stimulation frequency allows muscles to contract with a force that is precisely tailored to the needs of the moment, whether lifting a feather or a heavy box.

Wave summation

At low firing rates, motor neurons induce individual twitch contractions in muscle fibers. These twitches occur when the muscle fibers have time to fully relax between stimuli, resulting in sporadic, gentle contractions. If the neuron fires before the muscle fibers fully relax after the first twitch, the second contraction is added to the previous one, resulting in wave summation and a more robust overall contraction. The resultant phenomenon is called temporal summation or wave summation, wherein a more powerful second contraction occurs if a second stimulus is received while the relaxation phase continues.

Incomplete tetanus

The duration of a single twitch determines the maximum time available for wave summation. For example, if a twitch lasts for 20 milliseconds, subsequent stimuli must be separated by less than 20 milliseconds, which means a stimulation rate of more than 50 stimuli per second. This rate is usually expressed in terms of stimulus frequency, the number of stimuli per unit of time. As the frequency of stimulation increases, the muscle enters a state known as incomplete tetanus. During incomplete tetanus, the muscle fibers only partially relax between twitches, leading to a fluttering yet more forceful contraction. This state is often sufficient for many everyday tasks that require moderate strength.

Complete tetanus

At even higher frequencies, the muscle reaches complete tetanus, where individual twitches become indistinguishable, and the muscle fiber is in a state of continuous, maximal contraction. This occurs at frequencies of around 80 to 100 times per second and is typically utilized in situations requiring maximum force. In complete tetanus, the muscle does not relax at all between stimuli, resulting in a smooth and sustained contraction plateau that represents the peak force a muscle can generate.

Treppe

Stronger, more sustained contractions can also occur through the Treppe or the staircase effect. This occurs when consecutive stimuli are provided at a frequency that allows complete relaxation between twitches, yet each subsequent twitch produces a slightly greater force. This effect is attributed to a gradual increase in calcium ions in the sarcoplasm and increased muscle enzyme efficiency.

 Core: Anatomy and Physiology

Muscles of the Anterior Neck

JoVE 14869

The anterior neck muscles are the group of muscles covering the front part of the neck. These muscles are classified into three subgroups. The first one is the superficial muscles, the most visible muscles in the front of the neck. It includes the platysma and sternocleidomastoid. The second group is the suprahyoid muscles, located above the hyoid bone. This group comprises the digastric, mylohyoid, geniohyoid, and stylohyoid. Lastly, the infrahyoid muscles are found below the hyoid bone and consist of the sternohyoid, omohyoid, sternothyroid, and thyrohyoid.

Superficial Muscles

The platysma is a broad, thin muscle stretching from the chest and shoulders to the lower jaw and face. It is located just under the skin, and it is one of the superficial muscles that form the expression of the neck. The sternocleidomastoid muscle, often abbreviated as SCM, is a prominent and easily visible muscle that runs along each side of the neck. As its name suggests, its attachment points are the sternum, clavicle, and mastoid process of the skull. This muscle is primarily responsible for rotating and flexing the head.

Suprahyoid Muscles

All the four suprahyoid muscles connect to the hyoid bone. The digastric muscle has two parts, one starting at the jaw and the other at the temporal bone, which joins at the intermediate tendon. The flat and triangular mylohyoid and the geniohyoid muscles start at the jaw. The stylohyoid muscle is paired and begins at the styloid process of the temporal bone.

Infrahyoid Muscles

The infrahyoid muscles comprise the sternohyoid, omohyoid, sternothyroid, and thyrohyoid. The sternohyoid muscle starts at the manubrium and the clavicle. The omohyoid muscle has two parts — the inferior belly originates near the suprascapular notch, inclining upwards and connecting with the superior belly via the central tendon. The thyrohyoid muscle begins at the thyroid cartilage and links to the hyoid bone alongside the sternohyoid and omohyoid muscles. Lastly, the sternothyroid muscle starts at the manubrium and first costal cartilage and connects to the thyroid cartilage.

 Core: Anatomy and Physiology

Neurons: The Cell Body and the Dendrites

JoVE 14885

A typical nerve cell comprises three main components: the cell body, dendrites, and the axon. The cell body, also known as the soma or perikaryon, serves as the central biosynthetic hub housing a nucleus surrounded by cytoplasm containing organelles commonly found in most cells. Notably, Nissl bodies, clusters of the rough endoplasmic reticulum and free ribosomes responsible for protein synthesis, are distinctive features of the neuronal cell body. As neurons age, aggregates of a brown pigment called lipofuscin, which is a product of lysosomes, accumulate. The cytoskeleton of the neurons consists of networks of intermediate filaments and microtubules called neurofilaments and microtubules, respectively. Microtubules are important for the transport of materials between the cell body and the axon. Neurofilaments form bundles called neurofibrils that provide structural support.

Dendrites, the short and branching extensions from the cell body, function to receive signals from other neurons, transforming them into short-distance electrical impulses known as graded potentials. Adorned with dendritic spines, tiny protrusions enhancing surface area for signal reception, dendrites have a plasma membrane rich in receptors for neurotransmitters—chemical messengers released by neighboring neurons upon stimulation.

 Core: Anatomy and Physiology

Disorders of the Nervous Tissue

JoVE 14902

Nervous tissue is a vital component of the human body's communication system, enabling us to perceive and respond to stimuli. However, like all other tissues, it is vulnerable to disorders and diseases that can significantly impact our neurological functioning.

Homeostatic Imbalances:

Alzheimer's disease manifests as a gradual decline in memory and cognitive abilities, attributed to the buildup of amyloid plaques and neurofibrillary tangles in the brain.

Parkinson's disease arises from the degeneration of dopaminergic neurons that play a vital role in motor function. This degeneration leads to tremors, rigidity, and challenges in initiating voluntary movements.

Multiple sclerosis is a condition where the immune system mistakenly attacks and damages the myelin sheath, a protective covering of nerve fibers in the brain and spinal cord. This process, known as demyelination, leads to compromised communication between nerve cells.

Infections:

Meningitis is a disease that results from an infection, causing the meninges - the protective layers around the brain and spinal cord - to become inflamed. This often leads to symptoms like a high temperature, headaches, and a stiff neck.

Poliomyelitis, a viral infection, targets the motor neurons in the spinal cord, resulting in muscle weakness and paralysis.

Toxicity:

Lead poisoning: High levels of lead exposure can cause brain and nervous system damage, leading to developmental delays, seizures, and behavioral issues.

Carbon monoxide poisoning: Carbon monoxide is a poisonous gas that can affect the brain and nervous systems, leading to symptoms such as headache, dizziness, and confusion.

Genetics:

Huntington's disease: an inherited disorder that causes the destruction of acetylcholine and GABA (gamma-aminobutyric acid) producing neurons, leading to involuntary movements, cognitive problems, and behavioral changes.

Tourette Syndrome: a genetic disorder that results in repetitive, involuntary movements and vocalizations called tics.

 Core: Anatomy and Physiology

Spinal Cord: Cross-sectional Anatomy

JoVE 14918

The cross-sectional anatomy of the spinal cord offers a detailed view of its complex structure and function within the central nervous system. At the core of the spinal cord lies the gray matter, characterized by its butterfly or "H"-shaped appearance in cross-section. This central region is enveloped by white matter, with the overall structure divided into symmetrical halves by the dorsal median sulcus and the ventral median fissure.

Gray Matter and its Components

Central to the gray matter is the gray commissure, encircling the central canal filled with cerebrospinal fluid (CSF). This commissure acts as a bridge connecting the left and right sides of the spinal cord. Projections from the gray commissure extend outward, forming the dorsal and ventral horns.

The dorsal horns contain multipolar interneurons and the axons of sensory neurons, which organize into dorsal rootlets and dorsal roots. These structures receive input from descending pathways and are pivotal in processing incoming sensory information. Dysfunction of the dorsal horns can lead to sensory abnormalities and contribute to conditions such as chronic pain syndromes.

The ventral horns comprise multipolar somatic motor neurons and interneurons essential for initiating motor responses. Their axons bundle into ventral rootlets and ventral roots, forming spinal nerves that transmit motor signals from the spinal cord to muscles, controlling voluntary movements and motor functions. The dorsal and ventral roots converge on each side of the spinal cord to form the spinal nerves. This union is a critical juncture where sensory inputs are integrated with motor outputs, facilitating the spinal cord's primary role in relaying information.

The gray matter exhibits lateral horns in certain regions of the spinal cord, notably the thoracic and upper lumbar areas. These contain autonomic motor neurons that give rise to preganglionic sympathetic neurons. These neurons are then projected to sympathetic ganglia located outside the spinal cord. As a result, these neurons play a crucial role in controlling various involuntary functions, such as heart rate, blood pressure, and respiration.

The quantity of gray matter within the spinal cord varies along its length, reflecting the density of neural connections necessary to innervate different body parts. Regions controlling the limbs, for example, contain a higher concentration of gray matter due to the complexity of motor and sensory functions required.

White Matter Organization

The white matter surrounds the gray matter and is segmented into dorsal, ventral, and lateral columns, known as funiculi. These columns house ascending and descending nerve fiber tracts dedicated to sensory and motor functions.

The ascending pathways carry sensory information from the peripheral nerves to the brain. These tracts process sensory data such as touch, temperature, pain, and proprioception. Notable ascending tracts include the dorsal columns fasciculus gracilis and fasciculus cuneatus.

Conversely, the descending tracts are involved in transmitting motor commands from the brain to the spinal cord. These tracts influence muscle movements and coordinate activities such as walking, grasping, and posture. Key descending tracts include the corticospinal tract.

In addition to containing portions of ascending and descending tracts, the lateral funiculus also includes important pathways like the lateral corticospinal tract and the rubrospinal tract, which play roles in voluntary movement control and modulation of limb movements.

 Core: Anatomy and Physiology

Hierarchy of Motor Control

JoVE 14936

The hierarchy of motor control refers to the different levels of organization and processing involved in controlling movement in the body. These levels range from higher cortical areas involved in planning and decision-making to lower spinal cord reflexes that respond automatically to external stimuli.

  1. Segmental Level: This is the lowest level of motor control and consists of spinal cord segments. It is responsible for generating simple reflexes and rhythmic movements such as walking. Neural circuits in the spinal cord, known as central pattern generators or CPGs, play an important role in generating and coordinating these rhythmic movements.
  2. Projection Level: This level is responsible for the execution of voluntary movements and comprises the primary motor cortex and the brainstem. These structures receive input from the segmental level and send output down to the spinal cord to initiate movement. It directly controls the movements through the segmental level and sends a copy of the information as internal feedback to the higher motor centers for planning and decision-making.
  3. Precommand Level: This is the highest level of motor control and is involved in the planning and initiation of complex movements. The precommand level includes the basal nuclei and the cerebrum, which integrates sensory information and past experiences to plan and execute movements.

Overall, the three levels of motor control — segmental, projection, and precommand — illustrate the different levels of organization and processing involved in controlling movement, from basic reflexes to complex, voluntary movements.

 Core: Anatomy and Physiology

Parasympathetic Division of the ANS

JoVE 14953

The parasympathetic division of the autonomic nervous system (ANS) regulates rest and digestion functions in the body. It works in opposition to the sympathetic division, promoting relaxation, conservation of energy, and digestion. The parasympathetic division consists of preganglionic fibers originating from specific cranial nerves (III, VII, IX, X) and the sacral spinal nerves (S2-S4). These fibers synapse with postganglionic neurons in the terminal ganglia, innervating various organs and tissues.

The sacral part of the parasympathetic division plays a crucial role in regulating the pelvic organs' visceral functions and the large intestine's distal half It is characterized by preganglionic axons originating from the anterior roots of the second through fourth sacral spinal nerves. As these preganglionic axons travel through the sacral spinal nerves, they branch off to form pelvic splanchnic nerves. These nerves synapse with parasympathetic postganglionic neurons in the terminal ganglia present in the walls of the innervated viscera. From these terminal ganglia, parasympathetic postganglionic axons innervate the smooth muscle and glands in the walls of the colon, ureters, urinary bladder, and reproductive organs.

The parasympathetic division influences functions such as:

  1. Pupil Constriction: Parasympathetic stimulation causes pupil constriction (miosis), allowing less light to enter the eye and improving near vision.
  2. Salivation: Parasympathetic innervation stimulates salivary glands, promoting the production and release of saliva, which aids in digestion and facilitates swallowing.
  3. Digestion: The parasympathetic division enhances digestive processes by increasing intestinal motility and promoting the secretion of digestive enzymes and gastric juices, facilitating the breakdown of food.
  4. Slowing of the Heart Rate: Parasympathetic stimulation slows the heart rate by decreasing the rate of electrical impulses generated by the sinoatrial (SA) node, promoting relaxation.
  5. Stimulation of Glandular Secretions: Parasympathetic activity stimulates the secretion of various glands, including lacrimal glands (tear production), nasal glands (mucus production), and digestive glands (gastric, pancreatic, and intestinal secretions).
  6. Emptying of the Bladder: Parasympathetic fibers stimulate the detrusor muscle in the bladder, causing its contraction and promoting bladder emptying during urination.

 Core: Anatomy and Physiology

An Overview of the Endocrine System

JoVE 14972

The endocrine system, a complex network of glands, orchestrates physiological balance within the body through the production and secretion of hormones. These hormones are chemical messengers in intercellular communication, acting as conduits between the secretory cells and distant target sites. They traverse the circulatory system by being released into the extracellular fluid, and their impact is specific to cells possessing receptors for a particular hormone.

The endocrine system collaborates with the nervous system to uphold homeostasis, regulating fundamental biological processes like metabolism, reproduction, and development. Although both systems engage in intercellular communication, they differ in their temporal dynamics and modes of information transmission. Nerve impulses rapidly propagate through the body via electrical signals and neurotransmitter release. Still, their effects promptly cease upon stimulus removal, as exemplified by the sensation of heat dissipating after removing a hand from a hot object.

Conversely, hormones exhibit diverse temporal patterns, with some triggering rapid responses, like the immediate surge of adrenaline in fight-or-flight situations. In contrast, others, such as reproductive hormones, instigate prolonged signaling, exerting influence even after removal. This intricate interplay between the endocrine and nervous systems underscores the sophistication of the body's regulatory mechanisms, ensuring adaptability and responsiveness to various physiological challenges.

 Core: Anatomy and Physiology

Adrenal Gland Disorders

JoVE 14988

Adrenal gland disorders manifest when the production of adrenal hormones deviates from the norm, resulting in either excessive or insufficient concentrations.

Adrenal insufficiency, characterized by insufficient cortisol and aldosterone production, leads to conditions like Addison's disease. This disorder, affecting the adrenal cortex, exhibits symptoms such as skin bronzing, dehydration, low blood pressure, fatigue, and weight loss. Congenital adrenal hyperplasia, a genetic ailment causing inadequate cortisol production, triggers premature puberty, enlarged penises in males, and ambiguous genitalia in females.

Conversely, excessive cortisol production causes Cushing's syndrome, characterized by rapid protein loss in muscles and bones, along with symptoms like diabetes, high blood pressure, excessive weight gain, and edema.

Another adrenal disorder arises from medullary chromaffin cell tumors known as pheochromocytoma, leading to elevated catecholamine levels. This condition prompts uncontrolled sympathetic nervous system activity, presenting symptoms like heightened blood pressure, excessive sweating, rapid heart rate, and anxiety.

Treatment approaches for adrenal disorders encompass a combination of medication, surgical interventions, and lifestyle modifications. The aim is to restore hormonal balance and alleviate symptoms, ensuring optimal functioning of the adrenal glands and overall well-being for affected individuals.

 Core: Anatomy and Physiology

Mesh Analysis with Current Sources

JoVE 15050

Mesh analysis becomes simpler when analyzing circuits with current sources, whether independent or dependent. The presence of current sources reduces the number of equations required for analysis. Two cases illustrate this:

Current Source in One Mesh: The analysis process is straightforward when a current source is found in only one mesh within the circuit. Mesh currents are assigned as usual, with the mesh containing the current source excluded from the analysis. Kirchhoff's voltage law (KVL) is applied to the remaining mesh, resulting in a linear equation. Since the current in the mesh with the source is equal in magnitude but opposite in direction, it allows for easy determination of the current in the first mesh.

Current Source Between Two Meshes: In cases where a current source lies between two meshes, the analysis can be simplified by creating a supermesh. This involves excluding the current source and any elements connected in series with it. Applying KVL to the supermesh yields a linear equation. Additionally, Kirchhoff's current law (KCL) is applied to a node where the branch with the current source is connected, providing another linear equation linking the two branch currents. Solving these equations provides the values of the mesh currents.

Critical properties of a supermesh include:

  • • The current source within the supermesh imposes a constraint equation necessary for solving the mesh currents.
  • • A supermesh does not have its own current; it encompasses currents from the individual meshes it encloses.
  • • KVL and KCL are applied to a supermesh like any other mesh.

 Core: Electrical Engineering

Current Dividers

JoVE 15073

In parallel electrical connections, resistors are linked between the same pair of nodes, creating an equal voltage across each resistor. Kirchhoff's current law is applied to these connections, establishing that the sum of currents through these resistors equals the source current. Utilizing Ohm's law, the source current is determined as the product of the source voltage and the sum of the reciprocals of individual resistances. This relationship simplifies the process of finding the current flowing through each resistor.

In a parallel arrangement, the source current distributes itself among the resistors in inverse proportion to their resistances, exemplifying the "current division" principle within a current divider circuit.

Parallel resistors can be approached as a single equivalent resistor, simplifying circuit complexity. The equivalent resistor's value is calculated as the product of individual resistances divided by their sum. Two parallel resistors yield an equivalent resistance by multiplying their resistances and dividing by their sum. This concept extends to circuits featuring more than two resistors in parallel, with the equivalent resistance determined using a similar approach.

Equation1

Conductance, the reciprocal of resistance, is a valuable parameter in series and parallel connections. In series resistors, the equivalent conductance stems from the product of individual conductances divided by their sum. Conversely, in parallel configurations, it represents the sum of the individual conductances. The use of conductance calculations offers convenience when handling parallel resistors. The equivalent conductance of parallel resistors becomes the sum of their individual conductances, mirroring the equivalent resistance calculation for series resistors.

In practical scenarios, the combination of resistors in both series and parallel configurations aids in simplifying complex networks. This simplification facilitates the analysis and design of electrical circuits while preserving the original network's current-voltage (i-v) characteristics within the simplified setup.

 Core: Electrical Engineering

Energy Stored in Inductors

JoVE 15089

An inductor is ingeniously crafted to accumulate energy within its magnetic field. This field is a direct result of the current that meanders through its coiled structure. When this current maintains a steady state, there is no detectable voltage across the inductor, prompting it to mimic the behavior of a short circuit when faced with direct current.

In terms of gauging the energy stored within an inductor, it is equivalent to the integral of the power delivered at every individual moment, all accumulated over a specific duration of time. Mathematically, energy stored in an inductor is expressed as

Equation1

Where w is the energy stored in the inductor, L is the inductance and i is the current passing through the inductor.

Ideal inductors have a noteworthy characteristic - they do not dissipate energy. This trait allows the energy stored within them to be harnessed at a later point in time. However, this ideal scenario is slightly marred when dealing with non-ideal inductors.

Non-ideal inductors exhibit a phenomenon known as winding resistance. This resistance stems from the coils of the conductor and presents itself in series with the inductance. While this winding resistance has the potential to contribute to energy dissipation, it is typically so minuscule that it can be conveniently overlooked in practical applications.

Additionally, non-ideal inductors also display winding capacitance. This is due to the capacitive coupling that occurs between the conducting coils. However, this winding capacitance is usually so minute that it can be disregarded, except when dealing with high frequencies.

 Core: Electrical Engineering

Phasor Relationships for Circuit Elements

JoVE 15107

Phasor representation is a powerful tool used to transform the voltage-current relationship for resistors, inductors, and capacitors from the time domain to the frequency domain. This transformation simplifies the analysis of alternating current (AC) circuits.

In the time domain, Ohm's law provides a fundamental relation between the current flowing through a resistor and the voltage across it:

Equation1

where V is the voltage, I is the current, and R is the resistance. In phasor representation, this relationship holds true as well, with the voltage and current phasors being in phase and following Ohm's law.

For an inductor, the relationship between the voltage across it and the current flowing through it is given by the rate of change of current. The sinusoidal function representing this relationship can be converted into its phasor in polar format. When comparing the current and voltage phasors for an inductor, it can be observed that the current lags the voltage by 90 degrees. Using Euler's identity, a fundamental formula in complex analysis, the current-voltage relationship in the phasor domain can be obtained.

Similarly, when charging a capacitor, the current passing through it is determined by the rate of change of voltage across it. Again, the sinusoidal function representing this relationship can be converted into its phasor in polar form. In the case of a capacitor, the phasor representations indicate that the current leads the voltage by 90 degrees. The relationship between the current and voltage phasors for a capacitor can be obtained by using the time derivative of the voltage.

 Core: Electrical Engineering

Design Example: Capacitance Multiplier Circuit

JoVE 15176

In integrated circuit technology, a capacitance multiplier is often utilized to produce a larger capacitance value when a small physical capacitance falls short. This is achieved by a circuit that multiplies capacitance values by a factor of up to 1000, such that a 10-pF capacitor can replicate the performance of a 100-nF capacitor.

The circuit illustrated in Figure 1 below incorporates two op-amps, with the first operating as a voltage follower and the second acting as an inverting amplifier.

Figure1

Figure 1: Capacitance Multiplier

The voltage follower functions to isolate the capacitance created by the circuit from the loading incurred by the inverting amplifier. Since no current enters the op amp's input terminals, the feedback capacitor carries the input current.

By applying Kirchhoff's Current Law (KCL), a relation between input and output voltage with respect to resistances can be established, which can be further substituted into the current expression. Rearranging the expressions aids in determining the input impedance. By selecting appropriate resistance values, an effective capacitance can be generated between the input terminal and ground that is a multiple of the physical capacitance.

To prevent op-amps from saturating, the effective capacitance must be limited by the inverted output voltage. As the capacitance multiplication increases, the maximum allowable input voltage must decrease. Capacitance multiplier circuits such as this one provide an efficient solution for generating larger capacitances without increasing the physical capacitance.

 Core: Electrical Engineering

Diabetes Mellitus: Type 2 and Gestational

JoVE 16320

Type 2 diabetes, characterized by insulin resistance, arises when the insulin receptors on cells lose responsiveness to insulin, diminishing the cell's capacity to take up glucose, resulting in elevated blood glucose levels. To receive a diagnosis of Type 2 diabetes, a series of blood glucose tests are necessary to assess whether the blood glucose falls within normal parameters. If the result is out of the normal range, a patient may be diagnosed as prediabetic or diabetic, depending on the severity of the glucose elevation.

Type 2 diabetes, more common in adults, is associated with risk factors such as a sedentary lifestyle, obesity, poor diet, and a family history of the condition. Similar to Type 1 diabetes, individuals with Type 2 diabetes are susceptible to complications like heart disease, amputations, kidney failure, and blindness.

Management of Type 2 diabetes includes lifestyle modifications, encompassing dietary changes and regular exercise. Pharmacological management, such as oral antidiabetic medications, may be necessary, and in some cases, insulin therapy may be prescribed.

Gestational diabetes mellitus (GDM) occurs during pregnancy and can lead to complications such as preterm labor, fetal overgrowth, and preeclampsia. While GDM typically resolves post-childbirth, both mothers and their babies face a lifelong risk of developing Type 2 diabetes, cardiovascular diseases, and obesity.

Comprehensive management strategies, including early diagnosis, lifestyle interventions, and appropriate medical treatments, are crucial in mitigating the risks associated with Type 2 diabetes and its gestational counterpart, GDM.

 Core: Anatomy and Physiology

Zebrafish Reproduction and Development

JoVE 5151

The zebrafish (Danio rerio) has become a popular model for studying genetics and developmental biology. The transparency of these animals at early developmental stages permits the direct visualization of tissue morphogenesis at the cellular level. Furthermore, zebrafish are amenable to genetic manipulation, allowing researchers to determine the effect of gene expression on the development of a vertebrate with a high degree of genetic similarity to humans.

This video provides a brief overview of the major phases of zebrafish development, with particular focus on the first 24 hours post fertilization (hpf). The discussion begins with a zygote consisting of a single cell, or blastomere, atop a large ball of yolk. Cleavage of the blastomere is then shown to produce an embryo containing thousands of cells within a matter of hours. Next, the dramatic cellular movements known as epiboly and gastrulation are explained, revealing how they contribute to reshaping a mass of cells into a moving embryo with a beating heart in just 1 day. The presentation follows embryo development through the hatching phase, when they become swimming, feeding larvae. Important considerations for caring for larvae are incorporated, including a brief review of how fish are raised to adulthood in a dedicated facility known as the nursery. Finally, the video concludes with some common techniques utilized for studying embryo development, demonstrating how zebrafish are used to help us better understand human development and disease.

 Biology II

Calcium Imaging in Neurons

JoVE 5203

Calcium ions play an integral role in neuron function: They act as intracellular signals that can elicit responses such as altered gene expression and neurotransmitter release from synaptic vesicles. Within the cell, calcium concentration is highly dynamic due to the presence of pumps that selectively transport these ions in response to a variety of signals. Calcium imaging takes advantage of intracellular calcium flux to directly visualize calcium signaling in living neurons.

This video begins with an overview of the key reagents used for this technique, known as calcium indicator dyes. The discussion includes an introduction to the commonly used dye Fura-2 and some basic principles behind how both ratiometric and non-ratiometric calcium indicators work. Next, a typical calcium imaging experiment is presented, from preparing the cells and dye to capturing and analyzing the fluorescent images. Finally, several experimental applications of calcium imaging are provided, such as the study of neuronal network activity and sensory processing.

 Neuroscience

Gene Silencing with Morpholinos

JoVE 5326

Morpholino-mediated gene silencing is a common technique used to study roles of specific genes during development. Morpholinos inhibit gene expression by hybridizing to complementary mRNAs. Due to their unique chemistry, morpholinos are easy to produce and store, which makes them remarkably cost effective compared to other gene silencing methods.

This video reviews proper experimental design when using these oligonucleotides. Following that, an explanation of morpholino microinjection techniques in zebrafish and the analysis of resulting phenotypes will be discussed. Finally, we showcase examples of specific applications where morpholino technology is used to model developmental disorders or to study tissue regeneration.

 Developmental Biology

Performing 1D Thin Layer Chromatography

JoVE 5499

Source: Laboratory of Dr. Yuri Bolshan — University of Ontario Institute of Technology

Thin layer chromatography (TLC) is a chromatographic method used to separate mixtures of non-volatile compounds. A TLC plate consists of a thin layer of adsorbent material (the stationary phase) fixed to an appropriate solid support such as plastic, aluminum, or glass1. The sample(s) and reference compound(s) are dissolved in an appropriate solvent and applied near the bottom edge of the TLC plate in small spots. The TLC plate is developed by immersing the bottom edge in the developing solvent consisting of an appropriate mobile phase. Capillary action allows the mobile phase to move up the adsorbent layer. As the solvent moves up the TLC plate, it carries with it the components of each spot and separates them based on their physical interactions with the mobile and stationary phases.

 Organic Chemistry

RNA-Seq

JoVE 5548

Among different methods to evaluate gene expression, the high-throughput sequencing of RNA, or RNA-seq. is particularly attractive, as it can be performed and analyzed without relying on prior available genomic information. During RNA-seq, RNA isolated from samples of interest is used to generate a DNA library, which is then amplified and sequenced. Ultimately, RNA-seq can determine which genes are expressed, the levels of their expression, and the presence of any previously unknown transcripts.

Here, JoVE presents the basic principles behind RNA-seq. We then discuss the experimental and analytical steps of a general RNA-seq protocol. Finally, we examine how researchers are currently using RNA-seq, for example, to compare gene expression between different biological samples, or to characterize protein-RNA interactions.

 Genetics

An Introduction to Cell Motility and Migration

JoVE 5643

Cell motility and migration play important roles in both normal biology and in disease. On one hand, migration allows cells to generate complex tissues and organs during development, but on the other hand, the same mechanisms are used by tumor cells to move and spread in a process known as cancer metastasis. One of the primary cellular machineries that make cell movement possible is an intracellular network of myosin and actin molecules, together known as “actomyosin”, which creates a contractile force to pull a cell in different directions.

In this video, JoVE presents a historical overview of the field of cell migration, noting how early work on muscle contraction led to the discovery of the actomyosin apparatus. We then explore some of the questions researchers are still asking about cell motility, and review techniques used to study different aspects of this phenomenon. Finally, we look at how researchers are currently studying cell migration, for example, to better understand metastasis.

 Cell Biology

Density Gradient Ultracentrifugation

JoVE 5685

Density gradient ultracentrifugation is a common technique used to isolate and purify biomolecules and cell structures. This technique exploits the fact that, in suspension, particles that are more dense than the solvent will sediment, while those that are less dense will float. A high-speed ultracentrifuge is used to accelerate this process in order to separate biomolecules within a density gradient, which can be established by layering liquids of decreasing density in a centrifuge tube.

The video will cover the principles of density gradient ultracentrifugation, including a procedure that demonstrates sample preparation, creation of a sucrose gradient, ultracentrifugation, and collection of fractionated analytes. The applications section discusses isolation of multi-protein complexes, isolation of nucleic acid complexes, and separation using cesium chloride density gradients.

Density gradient ultracentrifugation is a common approach to isolate and purify cell structures for biochemical experiments. The technique uses a high-speed, or ultra, centrifuge to nondestructively separate cellular components in a density gradient. This video describes the principles of density gradient ultracentrifugation, provides a general procedure using a sucrose gradient, and discusses some applications.

Let's start by examining the principles of ultracentrifuges and density gradients. A suspension contains particles in a liquid solvent. Because of gravity, particles denser than the solvent sediment out while those less dense than the solvent float. The greater the difference in density between the particle and the solvent, the faster the separation.

An ultracentrifuge contains a unit called a rotor, which rotates at highly controlled speeds, simulating a strong gravitational field. Within this field, the differences in density between particles and the solvent are magnified.

The strength of the field depends on the speed of rotation. Even a small rotor at a relatively low rotational speed can create a force thousands of times stronger than the earth's gravitational field.

If a tube contains several liquids of different densities, centrifugation will keep them in separate layers in order of density, with the densest liquid closest to the base. Such a layering of multiple liquids is called a "density gradient." There are two types. In step gradients, liquids of decreasing density are carefully layered on top one another. In continuous gradients, the liquids are mixed in varying proportions, so the density decreases smoothly from the base upwards.

Cellular organelles can be separated using a step gradient, through "isopycnic density-gradient centrifugation." This is the simplest and most common centrifugation procedure.

This procedure is used to separate the cellular structures. The more dense the organelle, the further it descends-with mitochondria at the top and nucleic acids towards the bottom.

Now that you know the principles behind the technique, let's see it in the lab.

Before the procedure is started, the manufacturer's speed and density ratings should be noted, and the ultracentrifuge checked for corrosion. This procedure uses a swinging-bucket rotor.

First, the cellular material is prepared by homogenizing the cells, which nondestructively releases their organelles. The homogenate may be fractionated through preliminary low-speed centrifugation, to remove low-density components. Next, the sucrose solutions are prepared.

Sucrose is added in increasing amounts so each solution is more concentrated, and therefore denser, than the preceding one. The exact densities of the solutions will depend on the components to be separated, which vary between organisms. The solutions should have densities between those of the components to be separated, with the last solution denser than the densest component of the analyte. Techniques for separating components denser than sucrose, like nucleic acids, are described in the applications.

The sucrose gradient is now created in a clean centrifuge tube. A pipette is used to draw up the most concentrated sucrose solution. With the tube held upright, the pipette tip is placed high against the wall, and the liquid dispensed steadily down. It's important that the working area is kept free of vibrations and other disturbances.

After replacing the tip, the remaining solutions are added in order of decreasing density. They are dispensed carefully to form distinct layers and avoid mixing. Finally, about half a milliliter of the cellular sample is added atop the gradient, and the tube is weighed. This is used to balance the weight distribution, the next step of the process.

Centrifugation should begin as soon as possible. The tube is placed in the rotor, which is then balanced by placing blank solutions of equal weight in opposing slots. The rotor is placed in the ultracentrifuge and the system sealed. The temperature and rotation speed and time are set. Typical values are 4 °C with a force of over 100,000 x g for 16 h.

After centrifugation, the tube is withdrawn from the rotor, taking care to keep it upright and undisturbed. The different cellular components have fractionated into discrete bands between the solution layers. The fractions can be collected with a syringe. Alternately, the bottom of the tube can be punctured with a fine, sterilized needle and the outflow collected in sterile tubes. The cellular components have now been isolated. They can be stored at -80 °C.

Now that we've seen the basic procedure, let's look at some applications.

A typical application is the isolation of multi-protein complexes in plant cells. In this example, complexes responsible for cyclic electron flow are being isolated from the thylakoid, the site of the light reaction in photosynthesis. This procedure uses discrete solutions of 14 to 45% sucrose. Centrifugation occurs over 100,000 x g for 14 h at 4 °C.

Because nucleic acids are denser than sucrose, isopycnic centrifugation cannot separate them from organelles nondestructively.

A different technique, known as "rate-zonal centrifugation" is used. It separates organelles based on their sedimentation rates, which depend not only on their densities, but also on their conformations. A continuous gradient is used to separate the components based on this property.

The procedural steps are similar to those for isopycnic cases. In this example, RNA-ribosome complexes are isolated using a continuous gradient of 5% to 20%, centrifuged at 230,000 x g. Centrifugation is interrupted after a few hours to prevent co-precipitation.

Nucleic acid strands can be separated from each other on the basis of density.

This is because strands rich in guanine and cytosine are denser than those rich in adenine and thiamine. In this case, the gradient cannot be made of sucrose, because sucrose is less dense than nucleic acids. Instead, cesium chloride gradients, typically from 1.65 to 1.75 g/mL are used, as they have sufficient density and a low viscosity.

Here we see plankton DNA being purified using a continuous cesium chloride gradient. Centrifugation occurs at over 1,000,000 x g for 18 h under vacuum.

You've just watched JoVE's video on ultracentrifugation with a sucrose density gradient. You should now understand how a density gradient works, how to construct a step gradient, and how to load and operate an ultracentrifuge. Thanks for watching!

 Biochemistry

Overview of Tissue Engineering

JoVE 5785

Tissue engineering aims to create artificial tissue from biomaterials, specific cells, and growth factors. These engineered tissue constructs have far-reaching benefits, with possibilities for organ replacement and tissue repair.

This video introduces the field of tissue engineering and examines the components of engineered tissue. This video also outlines some prominent methods used to create the tissue scaffold, introduce a cell population, and encourage growth and proliferation. Finally, some key challenges and important applications of the technology are demonstrated.

 Bioengineering

Molecular Geometry and Dipole Moments

JoVE 11683

The VSEPR theory can be used to determine the electron pair geometries and molecular structures as follows:

  1. Write the Lewis structure of the molecule or polyatomic ion.
  2. Count the number of electron groups (lone pairs and bonds) around the central atom. A single, double, or triple bond counts as one region of electron density.
  3. Identify the electron-pair geometry based on the number of electron groups.
  4. Use the number of lone pairs to determine the molecular structure. If more than one arrangement of lone pairs and chemical bonds is possible, choose the one that will minimize repulsions.

Dipole Moment of a Molecule

When atoms with different electronegativities form a bond, the electrons are pulled toward the more electronegative atom, leaving one atom with a partial positive charge (δ+) and the other atom with a partial negative charge (δ–). Such bonds are called polar covalent bonds, and the separation of charge gives rise to a bond dipole moment. The magnitude of a bond dipole moment is represented by the Greek letter µ and is given by:

 μ = Qr

where Q is the magnitude of the partial charges (determined by the electronegativity difference), and r is the distance between them. Dipole moments are commonly expressed in debyes, where one debye is equal to 3.336 × 10−30 C·m.

The bond dipole moment is a vector represented by an arrow pointing along the bond from the less electronegative toward the more electronegative atom, with a small plus sign on the less electronegative end.

A whole molecule may also have a separation of charge, depending on its molecular structure and the polarity of each of its bonds. Such molecules are said to be polar. The dipole moment measures the extent of net charge separation in the molecule as a whole. In diatomic molecules, the bond dipole moment determines the molecular polarity.

When a molecule contains more than one bond, the geometry must be taken into account. If the bonds in a molecule are arranged such that the vector sum of their bond moments equals zero, then the molecule is nonpolar (e.g., CO2). The water molecule has a bent molecular structure, and the two bond moments do not cancel. Therefore, water is a polar molecule with a net dipole moment.

This text has been adapted from Openstax, Chemistry 2e, Section 7.6 Molecular Structure and Polarity.

 Core: Organic Chemistry

Chirality at Nitrogen, Phosphorus, and Sulfur

JoVE 11729

Chirality is most prevalent in carbon-based tetrahedral compounds, but this important facet of molecular symmetry extends to sp3-hybridized nitrogen, phosphorus and sulfur centers, including trivalent molecules with lone pairs. Here, the lone pair behaves as a functional group in addition to the other three substituents to form an analogous tetrahedral center that can be chiral.

A consequence of chirality is the need for enantiomeric resolution. While this is theoretically possible for all chiral amines, it is, in practice, difficult to separate the enantiomers of most chiral amines. This is due to pyramidal or nitrogen inversion, where the enantiomers are readily convertible from one form to another at room temperature, as the barrier to interconversion is ~25 kJ/mol. To briefly summarize the mechanism of this conversion, when the enantiomer passes through the transition state for inversion, the central nitrogen atom is sp2 hybridized, with its unshared electron pair occupying a p orbital. Therefore, ammonium salts that have no lone pair do not exhibit this phenomenon, and such quaternary chiral salts can be resolved into individual (relatively stable) enantiomers. Further, sp3 phosphorus and sulfur compounds, despite their lone pair, possess a high barrier for interconversion. Hence, their enantiomeric resolution is feasible.

Recall that enantiomers are non-superposable and are, therefore, different compounds with distinct identities. The nomenclature of chiral nitrogen, phosphorus and sulfur centers is like that of chiral carbon centers. The process of naming their enantiomers follows the Cahn–Ingold–Prelog rules or (R-S system), which involves three steps. The three steps are the same as with carbon centers—namely,  assignment of priorities to the substituent groups, the orientation of the lowest-priority substituent away from the observer, and determining whether the priority sequence of the other three groups at the chiral center is clockwise or counterclockwise. However, in chiral centers with a lone pair, the lone pair is always assigned the lowest priority, as compared to hydrogen in systems without a lone pair. Accordingly, the molecule is rotated such that the lone pair points away. As with carbon, the chiral center is the R configuration if the one-two-three sequence is clockwise and the S configuration if the sequence is counterclockwise.

 Core: Organic Chemistry

Introduction to Electrophilic Addition Reactions of Alkenes

JoVE 11769

The double bond in a simple, unconjugated alkene is a region of high electron density that can act as a weak base or a nucleophile. The filled π orbital (HOMO) of the double bond can interact with the empty LUMO of an electrophile. A bonding interaction occurs when the electrophile attacks between the two carbons; the electrophile then accepts a pair of electrons from the π bond and undergoes addition across the double bond, yielding a single product.

Addition and elimination reactions can be considered to exist in a temperature-dependent equilibrium, which can be better understood from the change in Gibbs free energy (ΔG) of the reaction. In addition reactions, one π bond is broken, and two σ bonds are formed. These reactions are usually exothermic because σ bonds are stronger than π bonds; thus, the enthalpy term (ΔH) is negative. The entropy term (−TΔS) is always positive: the number of molecules decreases, leading to a negative ΔS, and T is always positive on the kelvin scale, so the negative of that product is a positive term overall. Consequently, the value of ΔG is dependent on the temperature of the system, and addition reactions are favored at low temperatures.

When an alkene undergoes halogenation, bonds are formed between carbon and the more electronegative halogens; thus, the carbon atoms are oxidized. Dihydroxylation, halohydrin formation, and epoxidation are also oxidation reactions. Conversely, the addition of hydrogen across the double bond in alkenes is a reduction reaction that yields the corresponding alkanes. In hydration and hydrohalogenation reactions, one of the carbon atoms is oxidized while the other is reduced; as a result, they are not classified as oxidation or reduction reactions. In the hydrobromination of but-2-ene, the acidic proton in HBr accepts a pair of electrons from the π bond. The proton is transferred to one of the carbons in the double bond, while the other carbon acquires a positive charge, resulting in a secondary carbocation intermediate. The bromide ion then reacts with the positive center to yield a racemic mixture of 2-bromobutane.

 Core: Organic Chemistry

Introduction to the Cytoskeleton

JoVE 11788

Overview of the Cytoskeleton

The cytoskeleton is a network of protein filaments present within the cell, having three distinct filaments ̶   microfilaments, microtubules, and intermediate filaments. Each has characteristic features that distinguish them, including the dynamics of their assembly and disassembly, mechanical properties, polarity, and the type of molecular motors associated with them. Earlier, they were thought to be present only in eukaryotic cells; however, their homologs were eventually found in prokaryotic cells. Studies on bacterial homologs of cytoskeletal proteins hypothesize that the cytoskeleton originated in bacteria and archaea.

Despite using the word 'skeleton,' the cytoskeleton is not a fixed structure. It is a dynamic and adaptive structure that participates in various cellular functions. These functions can broadly be categorized as i) Spatial organization of cellular content, ii) Connecting the cell to its external environment both physically and biochemically iii) Generation of coordinated forces that help in cell movement and change in cell shape. Although these filaments are organized into networks that resist deformation, they can undergo rapid reorganization in response to external signals or forces. 

Microfilaments or filamentous actin (F-actin) are right-handed spiral filaments of globular actin (G-actin) monomers. These are polar filaments owing to the rate of polymerization at either end. These filaments steadily elongate to produce a strong sustained force required to carry out motility and cell shape changes. Microtubules are hollow cylindrical structures having thirteen protofilaments made up of alpha-beta-tubulin heterodimers. Microtubules are known to have the most complex assembly and disassembly dynamics. Unlike microfilaments, the microtubules rapidly switch between polymerization and depolymerization. The microtubule dynamics are regulated by a structure known as Microtubule Organizing Centres (MTOCs). The third component, intermediate filaments, are long fibrous proteins composed of multiple subunits formed through multistep processes. These filaments are generally static structures, their dynamics regulated through post-translational modifications like phosphorylation and dephosphorylation.

 Core: Cell Biology

Microtubule Formation

JoVE 11908

Microtubules are dynamic structures that undergo continuous assembly and disassembly. They originate from specialized multi-protein complexes known as microtubule organizing centers or MTOCs. Within the MTOC, the point of origin of the microtubule is known as the minus end, while the end radiating outward is the plus end. Microtubules serve two primary functions — the organization of spindle complexes to separate sister chromatids during mitotic or meiotic cell division and the formation of locomotory appendages, like cilia and flagella.

MTOCs are found in both prokaryotic and eukaryotic organisms. However,  some lower eukaryotes, like most fungi, lack organized MTOCs. Instead, they have organized centrosomes consisting of centrioles and the pericentriolar material. In animal cells, the structure and location of MTOCs, vary within different cell types depending on the function of the microtubules.

Microtubule Nucleation

The nucleation of microtubules occurs within the MTOCs, i.e., the centrioles, where different γ-tubulin complex proteins interact with γ-tubulin subunits to form the γ-tubulin-ring complex (γ-TRC). Nucleation is initiated when the α-tubulin subunit of the αβ-tubulin heterodimer attaches to the γ-TRC. Several intrinsic and extrinsic factors influence microtubule nucleation. Intrinsic factors like the α- and β-tubulins isotype incorporated; the concentration of free αβ-tubulin heterodimers, the post-translational modifications, and the microtubule-associated proteins (MAPs) affect the microtubule nucleation dynamics. Extrinsic factors like temperature, pH, and microtubule interfering drugs are also responsible for the rate of microtubule polymerization or depolymerization.

 Core: Cell Biology

Preparation of Diols and Pinacol Rearrangement

JoVE 11928

Compounds bearing two hydroxyl groups are known as diols. When the hydroxyl groups are located on adjacent carbon atoms, the diols are called vicinal diols or glycols. Under acidic conditions, vicinal diols undergo a specific reaction called pinacol rearrangement.

The reaction begins with transferring a proton from the acid catalyst to one of the hydroxyl groups, producing an oxonium ion.

Figure1

In the second step, the oxonium ion loses H2O, forming a tertiary carbocation intermediate.

Figure2

In the following step, a methyl group migrates from one carbon to the adjacent carbon, producing a resonance-stabilized cation intermediate where carbon and oxygen have complete octets of valence electrons.

Figure3

Finally, the subsequent proton transfer from the resonance-stabilized cation intermediate to the solvent water completes the reaction to give pinacolone.

Figure4

Since the reaction causes an overall change in the carbon backbone, it is termed a rearrangement.

 Core: Organic Chemistry

Mechanism of Filopodia Formation

JoVE 12253

Filopodia are thin, actin-rich cellular protrusions that play an important role in many fundamental cellular functions. They vary in their occurrence, length, and positioning in different cell types, suggesting their diverse roles.

Their main function is to guide migrating cells during normal tissue morphogenesis or cancer metastasis by recognizing and making initial contacts with the extracellular matrix. However, they can also act as stationary cell anchors or help to establish communication between two cells. For example, the formation of filopodial bridges between two adjacent endothelial or epithelial cells is important for the subsequent establishment of adherens junctions. In addition, filopodia can also help the cells to reach and internalize distant targets, such as pathogens.

Filopodia Formation and Disassembly

Actin nucleation, elongation, and bundling are critical for filopodia formation and function. The filopodia core comprises 12-20 actin filaments spaced 12 nm apart. Bundling of these filaments by fascin is crucial to give filopodia the rigidity to maintain its elongated shape.

The disassembly of actin filaments during the retraction phase involves periodic helical and rotational motion of the actin shaft and is regulated by several factors. For instance, capping proteins promote filopodial retraction by shielding the polymerizing ends of the filaments from further elongation. RhoA kinase activity also regulates actin polymerization within filopodia.

 Core: Cell Biology

Diazonium Group Substitution with Halogens and Cyanide: Sandmeyer and Schiemann Reactions

JoVE 12539

Arenediazonium substitution reactions occur when the diazonium group is substituted by various functional groups such as halides, hydroxyl, nitrile, etc. For instance, arenediazonium salts react with copper(I) salts of chloride, bromide, or cyanide to form corresponding aryl chlorides, bromides, and nitriles. These reactions are named Sandmeyer reactions. Although the mechanism of this reaction is complicated, as illustrated in Figure 1, they are believed to progress via an aryl copper intermediate.

Figure1

Figure 1. The formation of aryl halides and nitriles via the Sandmeyer reaction

The hydrolysis of aryl nitriles into carboxylic acids is very facile. In this context, the Sandmeyer reaction plays a crucial role in converting aryl amines to substituted benzonitrile, as shown in Figure 2.

Figure2

Figure 2. The conversion of aryl amines to substituted benzonitrile.

Sandmeyer reactions are not used to prepare aryl fluorides and aryl iodides. For aryl fluorides, the arenediazonium salts react with fluoroboric acid to give diazonium fluoroborate in the form of a precipitated salt. As depicted in Figure 3, these residues are isolated, dried, and heated until they decompose to form the corresponding aryl fluorides. Such reactions are named Schiemann reactions. Similarly, arenediazonium salts react with potassium iodide to form aryl iodides.

Figure3

Figure 3. The formation of aryl fluoride via the Schiemann reaction.

 Core: Organic Chemistry

Average and Instantaneous Velocity Vectors

JoVE 12631

To calculate other physical quantities in kinematics, the time variable must be introduced. The time variable not only allows us to state where an object is (its position) during its motion, but also how fast it’s moving. The speed at which an object is moving is given by the rate at which the position changes with time. For each position, a particular time is assigned. If the details of the motion at each instant are not important, the rate is usually expressed as the average velocity v. This velocity vector is simply the total displacement between two points divided by the time taken to travel between them, also known as the elapsed time.

However, since objects in the real world move continuously through space and time, the velocity of an object at any single point is an important parameter. The quantity that tells us how fast an object is moving anywhere along its path is known as the instantaneous velocity, often simply called velocity. It is the average velocity between two points on a path, in the limit where the time (and therefore the displacement) between the two points approaches zero. Like average velocity, instantaneous velocity is a vector with a dimension of length per time.

This text is adapted from Openstax, University Physics Volume 1, Section 4.2: Acceleration Vector.

 Core: Physics

cAMP-dependent Protein Kinase Pathways

JoVE 13325

Cyclic Adenosine Monophosphate (cAMP) is an essential second messenger that activates protein kinase A (PKA) and regulates various biological processes. A single epinephrine molecule binds to GPCR and activates several heterotrimeric G proteins, each stimulating multiple adenylyl cyclase, amplifying the signal, and synthesizing large numbers of cAMP molecules. Small changes in cAMP concentration affect PKA activity. The binding of four cAMP molecules induces a conformational change in PKA, dissociating the catalytic subunits from the regulatory subunit. Activated PKA can now phosphorylate serine/threonine residues of downstream target proteins and stimulate them to produce an appropriate cellular response. PKA can generate distinct responses in different cells by activating specific target proteins, even when stimulated by the same extracellular ligand.

In liver and muscle cells, epinephrine-bound G protein-coupled receptors (GPCR) cause a rise in cAMP levels. The increased cAMP further activates PKA to promote glucose mobilization in two ways.

  1. It phosphorylates glycogen phosphorylase kinase (GPK) and activates it. GPK further phosphorylates and activates glycogen phosphorylase (GP), which catalyzes the breakdown of glycogen into glucose 1-phosphate.
  2. PKA also phosphorylates and inhibits glycogen synthase (GS) and prevents glycogen synthesis.

In addition, PKA phosphorylates an inhibitor of phosphoprotein phosphatase (IP). The phosphorylated IP  binds and blocks phosphoprotein phosphatase, preventing it from dephosphorylating GPK, GP, or GS.

Once the extracellular stimulus is removed, cAMP levels decrease, switching off PKA. Inactive PKA cannot activate phosphoprotein phosphatase inhibitors. Thus, phosphoprotein phosphatase becomes active and removes phosphates from enzymes involved in glycogen degradation and synthesis. The dephosphorylation promotes glycogen synthesis and prevents glucose mobilization.

Contrarily to liver and muscle cells, epinephrine-induced activation of PKA in adipose cells leads to phosphorylation and activation of the enzyme lipase. The activated enzyme breaks down stored triglycerides to produce free fatty acids, which are used as an energy source by the kidney, heart, and muscle cells.

 Core: Cell Biology

Plasmodesmata

JoVE 13369

In a multicellular organism, cells must communicate to work together in a coordinated manner. One way that cells communicate is through direct contact with other cells. The points of contact that connect adjacent cells are called intercellular junctions.

Intercellular junctions are a feature of fungal, plant, and animal cells. However, different types of junctions are found in different kinds of cells. Intercellular junctions found in animal cells include tight junctions, gap junctions, and desmosomes. The junctions connecting plant cells are called plasmodesmata (singular = plasmodesma). While they are functionally similar to animal gap junctions, they differ structurally.

Plasmodesmata are passageways that connect adjacent plant cells. Just as two rooms connected by a doorway share a wall, two plant cells connected by a plasmodesma share a cell wall.

The plasmodesma “doorway” creates a continuous network of cytoplasm—like air flowing between rooms. It is through this cytoplasmic network—called the symplast—that most nutrients and molecules are transferred among plant cells.

A single plant cell has thousands of plasmodesmata perforating its cell wall, making a giant communication network across the entire plant. The number and structure of plasmodesmata vary across cells and change in individual cells. Most water and nutrients that move through a plant are transported by vascular tissue—xylem and phloem. However, plasmodesmata also transport these materials among cells and ultimately throughout the plant.

Plasmodesmata are versatile and continuously alter their permeability. In addition to water and small molecules, they can transport specific macromolecules, such as receptor-like protein kinases, signaling molecules, transcription factors, and RNA-protein complexes.

As cells grow, the density of their plasmodesmata decreases unless they produce secondary plasmodesmata. Certain parasitic plants develop secondary plasmodesmata that connect them to hosts, allowing them to extract nutrients.

 Core: Cell Biology

X-ray Diffraction of Biological Samples

JoVE 13386

X-ray diffraction or XRD is an analytical tool that utilizes X-rays to study ordered structures such as crystalline organic and inorganic samples, polycrystalline materials, proteins, carbohydrates, and drugs.

According to Bragg's law, when X-rays strike the sample positioned on a stage, the rays are  scattered by the electron clouds around the sample atoms. The  X-ray diffraction or scattering is caused by constructive interference of the X-ray waves that reflect off the internal crystal planes. The scattering angle of the rays in specific directions by the atoms within the sample provides a diffraction pattern representing the electron density due to the atoms and bonds within the crystal. This information about the crystal lattice arrangement helps identify the material from its diffraction pattern, which resembles  a fingerprint. The material can be identified by comparing this fingerprint with a database, such as  the International Center for Diffraction Database (ICDD). XRD can look at the size, shape, and internal structure and help in the basic structural determination of a sample or unknown material.

For a single crystal, a typical diffraction pattern comprises spots that are 2D slices of 3-dimensional spheres. A computer program can integrate the resulting spots to determine the shape and intensity of the diffracted X-rays. Whereas, in a powder sample, the X-rays interact with many tiny crystals in random orientations. Hence, instead of spots, a circular diffraction pattern is observed. The intensities of the diffracted circles are plotted against the angles between the ring of the beam axis, giving a 2D plot known as a powder pattern. Amorphous or non-crystalline powders give a broad peak due to scattering from different directions compared to ordered scattering in crystalline powders, giving sharp concentric rings.

Certain precautions should be taken while studying biological samples using this technique. First of all, well-ordered crystals are a prerequisite for deciphering the protein structure through single-crystal XRD. Generally, single small crystals with well-defined faces, if they contain heavier atoms, will suffice and give a good diffraction pattern. For organic compounds, however, the crystals need to be larger. Without well-defined, viable crystals, this technique is not feasible. Some molecules are inherently more crystalline than others; thus, the difficulty of obtaining X-ray quality crystals can vary between compounds and is the main limiting factor of this technique. Further, the powder samples should be fine powder and have no clumps.

 Core: Cell Biology

Preparation of Samples for Electron Microscopy

JoVE 13402

To be visualized by an electron microscope, either transmission or scanning, biological samples need to be fixed (stabilized) so the electron beam does not destroy them and dried thoroughly (desiccated/dehydrated) so the vacuum does not affect them. Fixation needs to be done as quickly as possible because the sample properties will start changing as soon as it is removed from its natural environment. For example, in a tissue sample, the oxygen levels begin decreasing, causing an altered appearance of mitochondria.

Advances in TEM Sample Preparation

Cryofixation

In this technique, small cells and organelles can be frozen by plunge freezing, while larger cells and tissues can be frozen using high pressure. Plunge freezing is achieved by placing a small amount of sample on a cold EM grid and then rapidly plunging into liquid ethane filled in a reservoir and surrounded by liquid nitrogen. The rapid freezing ensures that water does not crystallize. In high-pressure freezing, liquid nitrogen is used to freeze the sample rapidly, and it is subjected to a ~2000 bar pressure to decelerate ice crystal formation. Cryofixed samples are sectioned in a cryo-ultramicrotome that consists of a frozen diamond knife.

Shadow casting

The rotating platform in a vacuum chamber holds the sample, which is then coated with heavy metal. The rotation of the sample ensures a metal coating on the front side and a shadow effect on the backside. The technique is particularly useful when studying ultrastructural details of bacteria or viruses, DNA, RNA, or isolated proteins.

Negative staining

Small samples, such as ribosomes, enzyme molecules, viruses, etc., can be imaged directly without sectioning. Samples are mixed with ammonium molybdate, uranyl acetate, or phosphotungstic acid, which are electron opaque. A metal layer is formed on the grid, except where the sample lies, creating a negative image of the sample. Contrast is achieved as a result of the pattern of metal distribution.

Preparing Desiccated and Non-desiccated Samples for SEM

SEM can image three types of samples, namely, (i) desiccated and conductive, (ii) desiccated and non-conductive, and (c) non-desiccated (wet). Naturally conductive samples that are desiccated can be easily mounted and imaged. However, desiccated and non-conductive samples must be coated with a conductive layer to protect the sample from overheating and improve the image quality. These samples are made conductive by coating the sample with a thin gold layer using a sputter coater. It uses argon ions to knock gold atoms off a gold plate and coat the sample surface with gold, making the surface conductive to electrons.

Wet samples can outgas in the vacuum environment of the microscope.

Wet SEM, also called Environmental SEM or ESEM, must control the water sublimation in the specimen during imaging. It is accomplished by first coating the sample with a thin layer of a membranous polymer that is transparent to electron beams and separates the wet specimen from the vacuum. Then the temperature is decreased, which decreases the saturated vapor pressure. These conditions allow instrument operation at a higher vacuum that preserves ultrastructural features of the specimen and reduces noise. A Deben Coolstage may also be used to image wet samples by freezing them to −25°C, which preserves the sample structure and prevents outgassing while electrons are bombarding it.

 Core: Cell Biology

What is Meiosis?

JoVE 13434

Meiosis is the process by which diploid cells divide to produce haploid daughter cells. In humans, each diploid cell contains 46 chromosomes, half from the mother and half from the father. Following meiosis, the resulting haploid eggs or sperm only contain 23 chromosomes; however, each of these chromosomes contains a unique combination of parental information that results from the meiotic process of crossing over.

Although meiosis shares similarities with mitosis—both rely on microtubules to partition chromosomes to opposite sides of a cell, which then divides to form a daughter cell pair—meiosis is only observed in the sex organs, while mitosis occurs in other tissue types of the body. In addition, the cells resulting from mitosis are genetically indistinguishable (save for random mutations) from their predecessor: crossing over does not occur, and all the daughter cells are diploid. In contrast, meiosis produces four cells that not only have half the number of chromosomes from their predecessor, but they also contain unique combinations of genetic material. No two meiotic products are identical, which helps account for the appearance and personality differences often seen between siblings in the same family.

 Core: Cell Biology

Renewal of Skin Epidermal Stem Cells

JoVE 13466

The skin is divided into epidermis, dermis, and hypodermis, the skin's outermost, middle, and inner layers. The human epidermal layer regularly undergoes renewal, where old, dead cells are replaced by new cells. Epidermal stem cells or EpiSCs divide and differentiate to restore the lost cells. For the renewal process, some EpiSCs continuously self-renew. In contrast, few others differentiate into transit-amplifying cells, which later form prickle or spinous cells, followed by granular cells, and finally, dead keratinized squames, which are routinely shed off the skin.

Many factors should be controlled and regulated to sustain the renewal process. These include stem cell and transit-amplifying or TA cell division rate, the time a cell takes to leave the basal layer and differentiate, and finally, the removal of old and dead cells. As the cells begin to differentiate, they express different keratin molecules. The basal cells express keratin markers K5 and K14, spinous cells have K1 and K10, granular cells have profilaggrin and loricrin, while the keratinized cells express filaggrin. In case of wound or tissue damage, the epithelial keratinocytes migrate to the wound site and re-epithelialize the site.

The epidermal stem cells are attached to the basal layer through β1-integrins expressed by stem cells. Integrins connect the actin filaments of the basal cells with the extracellular matrix.  In addition to attachment, the integrins also play a role in signaling pathways. Other signaling pathways include Wnt, Hedgehog, Notch, BMP, and EGF, which are required to renew and maintain epidermal cells. Mutation in any of these pathways might cause skin cancer.

 Core: Cell Biology

iPS Cell Differentiation

JoVE 13483

The ability of induced pluripotent stem cells or iPSCs to differentiate into most body cell types has stimulated repair and regenerative medicine research over the past few decades. iPSC-derived blood cells, hepatocytes, beta islet cells, cardiomyocytes, neurons, and other cell types can repair injuries or regenerate damaged tissue in diseases such as diabetes and neurodegenerative disorders.

iPSCs have been successfully used to treat age-related macular degeneration (AMD), a form of blindness. This disorder is caused by the loss of retinal pigment epithelium (RPE) due to aging. Skin cells from an AMD patient were isolated and reprogrammed to form iPSCs, which were then differentiated into RPE cells. When these newly formed RPE cells were transplanted into the patient's retina, they restored the patient's vision.

iPSCs have shown potential for treating sickle cell anemia using a patient's own cells. Researchers reprogrammed the bone marrow stromal cells of a sickle cell anemia patient to form iPSCs. The mutation that causes the sickle cell phenotype in iPSCs was corrected using the CRISPR-Cas system. These iPSCs, when differentiated into erythroid cells, expressed the normal β-globin protein.

iPSC-derived cells are also being explored for their potential in cancer therapy. Tissues destroyed in a cancer patient due to radiation or chemotherapy can be replaced using cells differentiated from the patient's cells. However, attempts to transplant such iPSC-derived cells have not to date been successful.

 Core: Cell Biology

Comparing Experimental Results: Student's t-Test

JoVE 14516

The t-test is a statistical method used to compare the sample mean with a population mean or compare two means from two data sets. The test statistic is calculated from the standard deviation, mean, and number of measurements in the data set at a selected confidence interval and then compared to a table of critical values at this confidence level. If the test statistic is smaller than the critical value, the null hypothesis is accepted. In this case, we state that the difference between the means is not statistically significant and therefore comes from indeterminate (random) errors. If the test statistic is higher than the critical value, the null hypothesis is rejected. In this case, we state that the difference between the means is statistically significant and cannot be explained by random errors. The difference comes from errors in the method, sampling, the analysts themselves, or true phenomenological differences. In statistics, t-tests can be performed on unpaired as well as paired data. Unpaired data are two sets of replicate measurements from the same source; paired data refer to data taken on the same samples or subjects from two different methods or at two different time points for comparison. When only one side of a normal distribution curve is used in the t-test, it is a one-tailed test. If both sides of the distribution are used in the t-test, the test is two-tailed.

 Core: Analytical Chemistry

Ladder Diagrams: Redox Equilibria

JoVE 14532

Ladder diagrams are useful tools for understanding redox equilibrium reactions, especially the effects of concentration changes on the electrochemical potential of the reaction. The vertical axis in the redox ladder diagrams represents the electrochemical potential, E. The area of predominance is demarcated using the Nernst equation.

Consider the Fe3+/Fe2+ half-reaction, which has a standard-state potential of +0.771 V. At potentials more positive than +0.771 V, Fe3+ predominates, whereas Fe2+ predominates at potentials more negative than +0.771 V. When the Fe3+/Fe2+ half-reaction is coupled with the Sn4+/Sn2+ reaction, the concentration of Fe3+ can be reduced by adding Sn2+ to excess. In this case, the potential of the resulting solution approaches +0.154 V down to +0.771 V, and Fe2+ and Sn4+ predominate.

To understand the interdependence between change in solution pH and electrochemical potential, consider the example of UO22+/U4+ half-reaction, whose electrochemical potential varies with the pH of the solution. As the pH of the solution decreases, the electrochemical potential increases, changing the dominant species from U4+ to UO22+.

 Core: Analytical Chemistry

Complexometric Titration: Overview

JoVE 14571

Complexometric titration involves the formation of a complex by reacting a metal ion with one or more ligands. A visual indicator often detects the end point of a complexometric titration. It is added to the metal solution before the titration, forming a stable metal–indicator complex and imparting color to the solution. As the titration approaches the equivalence point, the excess of the added ligand displaces the indicator from the metal–indicator complex, releasing the free indicator. The free indicator formation changes the color of the solution rapidly, signaling an endpoint. O,O′-dihydroxyazo compounds, such as eriochrome black T and calmagite, are frequently used as indicators in complexometric titrations.

 Core: Analytical Chemistry

Types of Coprecipitation

JoVE 14587

Coprecipitation is the contamination of a precipitate by otherwise soluble species and occurs via different processes. In colloidal precipitates, coprecipitation occurs via surface adsorption. For instance, barium sulfate has a primary layer of adsorbed barium ions and a secondary layer of nitrate counterions. This results in contamination of the precipitate by barium nitrate.

Sometimes, ions in a crystal lattice can undergo isomorphous replacement by inclusions of similar charge and size. For example, during the precipitation of cadmium sulfide, manganese in the solution can replace cadmium to form a mixed crystal. Similarly, magnesium ammonium phosphate can be contaminated by potassium ions in the solution, resulting in mixed crystals of magnesium potassium phosphate. The formation of mixed crystals can be prevented by removing the interfering ion or using a different precipitant.

Coprecipitation can also occur via occlusion, where foreign ions are trapped within the growing lattice. Similarly, pockets of solution are trapped between adjacent crystals in mechanical entrapment. Slow precipitation can minimize occlusion, whereas rapid dissolution and reprecipitation in a clean, fresh solvent can remove it.

 Core: Analytical Chemistry

Electrodeposition

JoVE 14620

Electrodeposition is a technique used to separate an analyte from interferents by electrochemical processes. Here, the analyte is a metal ion that can be deposited on an electrode immersed in the sample solution. The electrochemical setup consists of an anode and a cathode. When an electric current is applied to the setup, oxidation occurs at the anode. At the cathode, which consists of a large metal surface, metal ions undergo reduction and deposit onto the surface.

Electrodeposition can separate metals with significant differences in reduction potentials. For example, copper and nickel in the same solution can be separated by electrodeposition. In this process, the solution is first made acidic. When an electric current is applied, water molecules in the solution undergo oxidation, producing oxygen, and copper is deposited onto a platinum cathode. After weighing the cathode to determine the weight of the copper, the copper is removed from the cathode. The solution is then made basic, nickel is deposited onto the same platinum cathode, and its weight is determined by weighing the electrode.

Electrodeposition can also be used to separate ions that simultaneously deposit on different electrodes, such as copper(II) and lead(II). Lead(II) ions undergo oxidation in a nitric acid solution and deposit as lead dioxide on a platinum anode. Meanwhile, copper(II) ions undergo reduction and deposit as copper metal on a platinum gauze cathode. The quantitative results of electrodeposition depend on the potential and sizes of electrodes, the deposition time, and the stirring rate.

 Core: Analytical Chemistry

Isotonic and Isometric Muscle Contractions

JoVE 14849

Two primary types of muscle contractions are isotonic and isometric, each serving unique functions and involving distinct mechanisms. Both isotonic and isometric contractions are integral to the body's complex system of movement and stability. Isotonic exercises contribute significantly to functional strength and movement, while isometric contractions are crucial for maintaining posture and joint stability.

Isotonic contractions

Isotonic contractions occur when a muscle changes length while the tension remains constant, typically moving something heavy. These contractions are further divided into concentric and eccentric contractions. Concentric contractions involve the muscle shortening as it exerts force, such as when you lift a dumbbell. The muscle fibers shorten, pulling on tendons and moving parts of the body closer together. Eccentric contractions, on the other hand, occur when a muscle lengthens while still under tension, like when lowering the dumbbell back down. Though often overlooked, eccentric contractions are crucial for controlled movements and can significantly contribute to muscle strengthening and coordination. Isotonic movements are fundamental to everyday activities, including walking, running, and lifting objects. Exercise routines focussing on isotonic contractions enhance muscle mass, strength, and mobility.

Isometric contractions

Isometric contractions, in contrast, occur when muscle length remains unchanged while muscle tension increases. During isometric exercises, the muscle does not noticeably change length, and the affected joint doesn't move, yet the muscle is still working. This type of contraction is exemplified by activities such as holding a plank or carrying an object in a steady position. Isometric training is particularly beneficial for stabilizing muscles and joints, improving postural support, and can be used in rehabilitation settings to maintain muscle strength without placing stress on injured or vulnerable joints.

Both isotonic and isometric contractions are integral to the body's complex system of movement and stability. Isotonic exercises contribute significantly to functional strength and movement, while isometric contractions are crucial for maintaining posture and joint stability.

 Core: Anatomy and Physiology

Muscles that Move the Head

JoVE 14870

The muscles that move the head are a dynamic and complex group of structures that work together to facilitate a wide range of head movements, including rotation, flexion, extension, and lateral bending.

The bilateral sternocleidomastoid, or SCM, and the suprahyoid and infrahyoid muscles are significant head flexors. The SCM muscles originate at the sternum and clavicle and attach to the mastoid process of the temporal bone. The SCM contracts bilaterally to bend the head forward, whereas unilateral contraction causes lateral flexion on the same side or head rotation to the opposite side.

The back of the neck is covered by muscles that connect the skull to the spinal column and pectoral girdle. These muscles can be categorized into three layers. The superficial layer includes the trapezius, splenius capitis, and splenius cervicis. The trapezius controls shoulder blade movements and assists in extending and tilting the head backward. The splenius capitis and splenius cervicis, originating from the lower cervical and upper thoracic vertebrae and inserting into the skull and upper cervical vertebrae, respectively, work together to extend the head and neck.

The deep layer includes the cervical transversospinalis muscles, such as the semispinalis capitis, semispinalis cervicis, and multifidus cervicis. These muscles primarily extend from the transverse processes of the vertebrae to the spinous processes of higher vertebrae or the occipital bone. The semispinalis capitis is particularly notable for its role in extending and rotating the head, contributing to the ability to look up and turn the head from side to side.

The deepest layer encompasses the suboccipital muscles, interspinales cervicis, and intertransversarii colli muscles. The suboccipital muscles, located just below the occipital bone, are crucial for fine motor control, aiding in minor adjustments of head position and proprioception — the sense of self-movement and body position.

 Core: Anatomy and Physiology

Neurons: The Axon

JoVE 14886

Axons are long, cytoplasmic processes of nerve cells capable of propagating electrical impulses known as action potentials. The cytoplasm or axoplasm of an axon contains neurofibrils, neurotubules, small vesicles, lysosomes, mitochondria, and various enzymes, all encased within the axolemma, the plasma membrane of the axon.

The axon attaches to the cell body at a cone-shaped elevation called the axon hillock. The initial part of the axon, closest to the hillock, is known as the initial segment. Nerve impulses often originate at the junction of the axon hillock and the initial segment, an area referred to as the trigger zone.

Axons can branch out, creating side branches or collaterals, allowing one neuron to communicate with multiple other cells. These branches end in fine extensions called telodendria or terminal branches, culminating at synaptic terminals for communication with other cells. This site of communication is known as a synapse.

Two transport systems facilitate the movement of substances between the cell body and the axon: slow axonal transport and fast axonal transport. Slow axonal transport moves materials from the cell body towards the axon terminals at a pace of 1–5 mm per day, providing new axoplasm for developing, regenerating, or mature axons.

On the other hand, fast axonal transport can move materials at a speed of 200–400 mm per day in both directions. It uses protein motors to move materials along the microtubules of the neuron's cytoskeleton. Anterograde (forward) fast axonal transport carries organelles and synaptic vesicles from the cell body to the axon terminals. Conversely, retrograde (backward) fast axonal transport brings materials from the axon terminals to the cell body for degradation or recycling.

 Core: Anatomy and Physiology

Anatomy of the Brain: Major Regions

JoVE 14903

The brain is the most complex organ in the human body. It consists of four main parts: the cerebrum, diencephalon, cerebellum, and brainstem.

The cerebrum is the largest section of the brain and divides into left and right hemispheres, separated by a deep fissure. The cerebral outer layer of grey matter — the cerebral cortex — comprises elevations called gyri and shallow groves called sulci. The inner portion of white matter includes long nerve fibers known as axons, which connect various areas of grey matter throughout both hemispheres.

The cerebellum is the second most prominent part of the brain, located at the back of the skull. An arbor vitae — a network of white fibers — connects its two hemispheres. The cerebellum separates into three lobes on each side — anterior, posterior & flocculonodular lobes. It aids in coordinating activities like fine-tuning skeletal muscle movements, equilibrium, and balance.

The diencephalon is a region that lies between the cerebrum and midbrain and houses vital structures such as the thalamus and hypothalamus. The thalamus acts as a relay station for sensory signals from all over the body to reach specific areas in the cortex. Concurrently, the hypothalamus helps regulate basic needs such as hunger or thirst by releasing hormones into other parts of the body.

The brainstem, located between the diencephalon and the spinal cord, consists mainly of the medulla oblongata, pons, and midbrain. It is responsible for vital functions such as respiration, heart rate, and digestion.

Protective brain coverings include the meninges and three membrane layers. Between the dura mater and arachnoid membranes is a space filled with cerebrospinal fluid, which helps cushion the brain from any mechanical forces and serves as a medium for exchanging nutrition and waste products.

The overall shape of the human brain can be divided into three parts: the cerebrum, cerebellum & brainstem. The average cranial capacity for male adults is about 1300 cubic centimeters (cc) and 1150 cc for females. White matter constitutes most of the volume, whereas 10-15% is grey matter with neurons that send signals throughout the body.

In conclusion, the human brain has multiple complex parts that are intertwined with one another to form a single unit responsible for various vital functions like breathing or movement coordination.

 Core: Anatomy and Physiology

Spinal Cord: Information Processing

JoVE 14919

The spinal cord is an integral hub for motor and sensory information that enables the brain to communicate with the peripheral nervous system (PNS). This communication consists of relaying sensory data and transmission of motor commands.

Sensory Information Processing

Sensory information processing begins at the sensory receptors located in the skin and other tissues, which detect somatic sensory stimuli such as touch, temperature, or pain. These receptors function as catalysts, initiating nerve impulses that are transmitted towards the spinal cord via sensory neurons. Upon reaching the spinal cord, these impulses enter through the dorsal root, which serves as the entry point for all sensory information from the peripheral body into the spinal cord. The impulses then proceed to the dorsal gray horn of the spinal cord, which acts as a processing and relay center for incoming sensory data. From the dorsal gray horn, the sensory impulse can follow one of several paths.

  • • Direct Ascension to the Brain: Impulses from stimuli such as fine touch or vibration enter the white matter of the spinal cord almost immediately upon arrival and ascend directly to the brain. This pathway allows the rapid relay of sensory information to the brain for processing and response.
  • • Interneuron Relay: Impulses from other stimuli, such as pain and temperature, may be transferred to interneurons located within the dorsal gray horn. The axons of these interneurons then project into the white matter, where they ascend to the brain. This path allows for the integration and modulation of sensory signals before they reach higher processing centers.
  • • Spinal Reflex Pathways: Another possible route for an impulse entering the dorsal gray horn is to engage the somatic motor neurons involved in spinal reflex pathways. This mechanism enables reflex actions, such as the withdrawal reflex from a painful stimulus, without direct brain intervention and providing an immediate response to certain sensory inputs.

Motor Information Processing

Motor output from the brain, which directs voluntary movements and reflex responses, descends into the spinal cord through white matter tracts. These descending motor signals convey instructions from the brain to various body parts.

The motor impulse travels down the spinal cord and enters the ventral gray horn, which houses motor neurons for transmitting signals to the muscles. From the ventral gray horn, the impulse moves into the ventral root, which serves as the exit point from the spinal cord to the spinal nerve. Finally, the impulse is carried through the spinal nerve to the target muscles, culminating in the intended response, whether it be muscle contraction, movement, or modulation of muscle tone.

 Core: Anatomy and Physiology

Direct Motor Pathways

JoVE 14937

The direct motor pathways, also known as the pyramidal tracts, are a group of neural pathways that originate in the brain and descend through the spinal cord. They control the voluntary movement of the body. There are two major direct motor pathways: the corticospinal and the corticobulbar tracts.

The corticospinal tract is responsible for the voluntary movement of the limbs and trunk. It originates in the cerebral cortex of the brain and descends through the cerebrum's internal capsule and the midbrain's cerebral peduncle. From there, it enters the spinal cord and travels down the length of the spinal column. As it descends, the corticospinal tract gives off branches to innervate different spinal cord levels. Most fibers end in the ventral horn, where they synapse with lower motor neurons that innervate the skeletal muscles.

The corticobulbar tract is responsible for the voluntary control of the face, tongue, and throat muscles. It originates in the cerebral cortex and descends through the internal capsule, like the corticospinal tract. However, instead of traveling down the spinal cord, it terminates in the brainstem, synapses with cranial nerves controlling the facial and head muscles.

Both the corticospinal and corticobulbar tracts are composed of upper motor neurons, which originate in the cerebral cortex, and lower motor neurons, which innervate the muscles.

 Core: Anatomy and Physiology

Cranial Part of Parasympathetic Division

JoVE 14954

The cranial part of the parasympathetic division plays a crucial role in regulating the visceral functions of the head and specific structures in the neck, thoracic, and abdominopelvic cavities. Preganglionic fibers of the parasympathetic division exit the brain through cranial nerves III (oculomotor), VII (facial), IX (glossopharyngeal), and X (vagus), delivering parasympathetic output to the respective visceral structures.

The vagus nerve (cranial nerve X) alone accounts for approximately 75 percent of all parasympathetic outflow from the central nervous system (CNS). It provides preganglionic parasympathetic innervation to structures in the neck, thoracic cavity, and abdominopelvic cavity, extending as far as the distal portion of the large intestine. The cranial parasympathetic ganglia associated with these cranial nerves are located close to the organs they innervate.

The ciliary ganglia are situated laterally to each optic (II) nerve. Preganglionic axons travel alongside the oculomotor (III) nerves to reach the ciliary ganglia. Postganglionic axons originating from these ganglia innervate the smooth muscle fibers within the eyeball.

The pterygopalatine ganglia are found laterally to the sphenopalatine foramen. They receive preganglionic axons from the facial (VII) nerve and send postganglionic axons to the nasal mucosa, palate, pharynx, and lacrimal glands.

The submandibular ganglia are located near the ducts of the submandibular salivary glands. They also receive preganglionic axons from the facial nerves and send postganglionic axons to the submandibular and sublingual salivary glands.

The otic ganglia are situated just inferior to each foramen ovale. Preganglionic axons from the glossopharyngeal (IX) nerves reach the otic ganglia, sending postganglionic axons to the parotid salivary glands.

 Core: Anatomy and Physiology

Structures of the Endocrine System

JoVE 14973

The intricate framework of the endocrine system encompasses a diverse array of glands, with their target tissues and organs strategically distributed throughout the body. Central to this network are the endocrine glands, specialized structures that lack ducts and release hormones directly into the interstitial fluid. Notably, the hypothalamus, a vital neuroendocrine organ situated in the brain, governs neural functions and serves as a potent source of hormonal regulation. Near the hypothalamus are the pituitary and pineal glands in the head region.

The thyroid and parathyroid glands in the neck region play pivotal roles in metabolism and calcium homeostasis, respectively. In the upper chest, the thymus influences immune system development. The adrenal glands are positioned above the kidneys. They contribute to stress response and electrolyte balance. The gonads, comprising the testes in males and ovaries in females, generate reproductive hormones.

In addition, the pancreas, nestled near the stomach, showcases versatility by functioning as both an endocrine gland regulating glucose levels and an exocrine gland aiding digestion. Beyond these recognized glands, cells within various organs and tissues—including the heart, adipose tissue, skin, and placenta—can release endocrine signals, expanding the system's influence across the body's diverse physiological domains. This intricate spatial distribution underscores the holistic nature of the endocrine system, working in concert to maintain balance and coordination throughout the organism.

 Core: Anatomy and Physiology

The Pineal Gland

JoVE 14989

The pineal gland, a diminutive endocrine structure named for its pinecone-shaped appearance, is situated atop the third ventricle within the diencephalon region of the forebrain. This gland, composed of secretory cells known as pinealocytes arranged in compact cords and clusters around dense particles of calcium salts, plays a pivotal role in hormonal regulation.

The primary secretion of the pineal gland is the hormone melatonin, derived from serotonin. The concentration of melatonin in the bloodstream exhibits a day-night pattern linked to light exposure. Notably, melatonin peaks at night when photoreceptors in the eye are devoid of light, inducing a sense of drowsiness.

Beyond its role in sleep regulation, melatonin exerts antioxidant effects, contributing to cellular protection. Additionally, this hormone significantly mediates mating behavior by influencing the timing and duration of the reproductive cycle, puberty, and sexual maturation. The intricate interplay between melatonin and various physiological processes underscores the multifaceted impact of the pineal gland on both circadian rhythms and reproductive functions.

 Core: Anatomy and Physiology

Source Transformation

JoVE 15051

Source transformation is a fundamental technique employed in circuit analysis, offering a valuable tool for simplifying complex electrical circuits. This technique involves the replacement of either a voltage source in series with a resistor by a current source in parallel with a resistor, or vice versa. The key concept here is that when the original sources are deactivated (turned off), the equivalent resistance at the circuit's end terminals remains the same.

It is essential to note that when performing source transformations, the direction of the current source arrow always points toward the positive terminal of the voltage source. This convention ensures consistency and aids in maintaining proper circuit orientation.

However, it is worth mentioning that source transformation is not applicable to ideal voltage sources, as they possess zero internal resistance. In contrast, nonideal voltage sources feature non-zero internal resistance, making them amenable to source transformation. Similarly, ideal current sources with infinite internal resistance cannot be substituted with finite voltage sources.

To illustrate the practical application of source transformation, consider a circuit connected to a non-ideal voltage source (Figure 1) and a non-ideal current source (Figure 2) individually. When the series resistance equals the parallel resistance, and the voltage across the voltage source adheres to Ohm's law, these non-ideal sources become equivalent to each other.

Figure1

Figure 1: Circuit with to non-ideal voltage source       

Figure2

Figure 2: Circuit with to non-ideal current source

Replacing the nonideal voltage source with the equivalent nonideal current source does not alter the voltage or current characteristics of any element within the circuit. This demonstrates the power of source transformation in simplifying circuit analysis without affecting overall circuit behavior.

 Core: Electrical Engineering

Equivalent Resistance

JoVE 15074

In circuit analysis, situations often arise where resistors are neither in series nor parallel configurations. To tackle such scenarios, three-terminal equivalent networks like the wye (Y) (Figure 1 (a)) or tee (T) and delta (Δ) (Figure 1 (b)) or pi (π) networks come into play. These networks offer versatile solutions and are frequently encountered in various applications, including three-phase electrical systems, electrical filters, and matching networks.

Figure1

The essence of these networks lies in their adaptability. They can be employed individually or integrated into more complex circuits. For instance, when dealing with a delta network but finding it more convenient to work with a wye network, a wye network can be superimposed onto the existing delta configuration. To determine the equivalent resistances within the wye network, ensuring that the resistance between each pair of nodes in the delta network equals the resistance between the same pair of nodes in the Y (or T) network is crucial.

The conversion process between these networks involves mathematical relationships that relate to the resistances. The transformation is facilitated by introducing an extra node 'n' and adhering to the conversion rule: In the Y network, each resistor is the product of the resistors in the two adjacent Δ (delta) branches, divided by the sum of the three Δ resistors. Conversely, to convert a wye network into an equivalent delta network, one can use the following conversion rule: Each resistor in the network is the sum of all possible products of Y resistors taken two at a time, divided by the opposite Y resistor.

These conversion formulas simplify for balanced networks where the resistances in both Y and Δ configurations are equal. This transformation doesn't involve adding or removing components from the circuit but substitutes mathematically equivalent three-terminal network patterns. It effectively transforms a circuit, allowing resistors to be analyzed as if they were either in series or parallel, facilitating the calculation of equivalent resistance if necessary.

 Core: Electrical Engineering

Series and Parallel Inductors

JoVE 15090

In electrical circuits, integrating inductors into the toolkit of passive elements requires navigating the intricacies of series and parallel combinations involving these components. Practical circuits often feature configurations of multiple inductors, and understanding how to determine their equivalent inductance is vital.

For a series connection of N inductors, each carrying the same current, applying Kirchhoff's voltage law unveils a crucial relationship. Substituting the expression for inductor voltages leads to an insightful conclusion: the equivalent inductance of this series arrangement is the simple summation of the individual inductances. This phenomenon mirrors how resistors combine in series, emphasizing a fundamental similarity between these passive elements.

Equation1

In a parallel connection of inductors, each bears an identical voltage; Kirchhoff's current law can be applied. This leads to a significant revelation: the reciprocal of the equivalent inductance is equal to the sum of the reciprocals of the individual inductances. This parallel combination rule mirrors the behavior of resistors in parallel, illustrating a harmonious approach to dealing with different passive elements.

Equation2

In practice, circuits often feature a blend of series and parallel inductors. When confronted with such scenarios, the equivalent inductance of the parallel inductors is calculated and added to the series inductance to obtain the overall equivalent inductance of the circuit.

 Core: Electrical Engineering

Kirchoff's Laws using Phasors

JoVE 15108

Analyzing AC circuits in electrical systems is a fundamental aspect of electrical engineering. In these circuits, AC power is supplied from a distribution panel and wired to various household appliances in parallel. To perform a comprehensive analysis, electrical engineers use Kirchhoff's voltage and current laws, which are equally applicable in AC circuits as in DC circuits.

Kirchhoff's voltage law (KVL) states that the sum of phasor voltages around a closed loop in an AC circuit equals zero. In the sinusoidal steady state, where AC voltages vary sinusoidally with time, these voltages can be represented in the time domain and then converted into phasor equivalents. This process ensures that the sum of the phasor voltages in a closed loop remains zero.

Equation1

Kirchhoff's current law (KCL) applies to circuit nodes, asserting that the total current entering a node equals the total current exiting the node. When expressed in phasor notation, the sum of phasor currents at a node also equals zero.

Equation2

These laws are essential tools for AC circuit analysis and enable engineers to work seamlessly in the frequency domain. They facilitate various tasks such as impedance combination, nodal and mesh analysis, superposition, and source transformation, which are crucial in designing and troubleshooting electrical circuits in residential and industrial settings.

 Core: Electrical Engineering

Degrees of Freedom

JoVE 17356

The degree of freedom for a particular statistical calculation is the number of values that are free to vary. As a result, the minimum number of independent numbers can specify a particular statistic. The degrees of freedom differ greatly depending on known and uncalculated statistical components.

For example, suppose there are three unknown numbers whose mean is 10; although we can freely assign values to the first and second numbers, the value of the last number can not be arbitrarily assigned. Since the first two numbers are independent with the third number being dependent, the dataset is said to have two degrees of freedom. In many statistical methods, the number of degrees of freedom is usually calculated as the sample size minus one. The degrees of freedom have broad applications in calculating standard deviation and statistical estimates in methods such as the Student t distribution and the Chi-Square distribution tests.

 Core: Analytical Chemistry

Zebrafish Maintenance and Husbandry

JoVE 5152

The zebrafish (Danio rerio) is a powerful vertebrate model system for studying development, modeling disease, and screening for novel therapeutics. Due to their small size, large numbers of zebrafish can be housed in the laboratory at low cost. Although zebrafish are relatively easy to maintain, special consideration must be given to both diet and water quality to in order to optimize fish health and reproductive success.

This video will provide an overview of zebrafish husbandry and maintenance in the lab. After a brief review of the natural zebrafish habitat, techniques essential to recreating this environment in the lab will be discussed, including key elements of fish facility water recirculation systems and the preparation of brine shrimp as part of the zebrafish diet. Additionally, the presentation will include information on how specific zebrafish strains are tracked in a laboratory setting, with specific reference to the collection of tail fin samples for DNA extraction and genotyping. Finally, experimental modifications of the zebrafish environment will be discussed as a means to further our understanding of these fish, and in turn, ourselves.

 Biology II

Rodent Stereotaxic Surgery

JoVE 5205

Stereotaxic (or stereotactic) surgery is a method used to manipulate the brain of living animals. This technique allows researchers to accurately target deep structures within the brain through the use of a stereotaxic atlas, which provides the 3D coordinates of each area with respect to anatomical landmarks on the skull. After the skull is exposed, anesthetized animals are mounted on a specialized instrument known as a stereotaxic frame, which enables the precise placement of experimental tools at the defined coordinates. Stereotaxic surgery is a versatile approach that can be used to generate lesions, manipulate gene expression, or deliver experimental agents to the brain.

This video-article provides a general overview of the principles behind stereotaxic surgery, including instructions for using a stereotaxic atlas and the stereotaxic frame, and an introduction to reading the Vernier scale for measurement of probe movements. The subsequent discussion outlines the steps required to perform the surgical procedure. Lastly, a broad range of technical applications are presented, such as the insertion of electrical probes to measure brain activity and genetic manipulation of brain tissue.

 Neuroscience

Genetic Engineering of Model Organisms

JoVE 5327

Transgenesis, or the use of genetic engineering to alter gene expression, is widely used in the field of developmental biology. Scientists use a number of approaches to alter the function of genes to understand their roles in developmental processes. This includes replacement of a gene with a nonfunctional copy, or adding a visualizable tag to a gene that allows the resultant fusion protein to be tracked throughout development.

In this video, the viewers will learn about the principles behind transgenesis, as well as the basic steps for introducing genetic constructs into an animal and targeting genes of interest. This is followed by the discussion of a protocol to create knockout mice. Lastly, some specific applications of transgenic technologies in the field of developmental biology will be reviewed.

 Developmental Biology

An Introduction to Learning and Memory

JoVE 5416

Learning is the process of acquiring new information and memory is the retention or storage of that information. Different types of learning, such as non-associative and associative learning, and different types of memory, such as long-term and short-term memory, have been associated with human behaviors. Studying these components in detail helps behavioral scientists understand the neural mechanisms behind these two complex phenomena.

JoVE's overview on learning and memory introduces common terminologies and a brief outline of concepts in this field. Then, key questions asked by behavioral scientists and prominent tools such as fear conditioning and fMRI are discussed. Finally, actual experiments dealing with aging, eradication of traumatic memories, and improvising learning are reviewed.

 Behavioral Science

Using a pH Meter

JoVE 5500

Source: Laboratory of Dr. Zhongqi He - United States Department of Agriculture

Acids and bases are substances capable of donating protons (H+) and hydroxide ions (OH-), respectively. They are two extremes that describe chemicals. Mixing acids and bases can cancel out or neutralize their extreme effects. A substance that is neither acidic nor basic is neutral. The values of proton concentration ([H+]) for most solutions are inconveniently small and difficult to compare so that a more practical quantity, pH, has been introduced. pH was originally defined as the decimal logarithm of the reciprocal of the molar concentration of protons Equation 1, but was updated to the decimal logarithm of the reciprocal of the hydrogen ion activity Equation 2. The former definition is now occasionally expressed as p[H]. The difference between p[H] and pH is quite small. It has been stated that pH = p[H] + 0.04. It is common practice to use the term 'pH' for both types of measurements.

The pH scale typically ranges from 0 to 14. For a 1 M solution of a strong acid, pH=0 and for a 1 M solution of a strong base, pH=14. Thus, measured pH values will lie mostly in the range 0 to 14, though values outside that range are entirely possible. Pure water is neutral with pH=7. A pH less than 7 is acidic, and a pH greater than 7 is basic. As the pH scale is logarithmic, pH is a dimensionless quantity. Each whole pH value below 7 is 10x more acidic than the next integer. For example, a pH of 4 is 10x more acidic than a pH of 5 and 100x (10 x 10) more acidic than a pH of 6. The same holds true for pH values above 7, each of which is 10x more basic (or alkaline) than the next lower whole value. For example, a pH of 10 is 10x more basic than a pH of 9.

 General Chemistry

An Overview of Epigenetics

JoVE 5549

Since the early days of genetics research, scientists have noted certain heritable phenotypic differences that are not due to differences in the nucleotide sequence of DNA. Current evidence suggests that these “epigenetic” phenomena might be controlled by a number of mechanisms, including the modification of DNA cytosine bases with methyl groups, the addition of various chemical groups to histone proteins, and the recruitment of protein factors to specific DNA sites via interactions with non-protein-coding RNAs.

In this video, JoVE presents the history of important discoveries in epigenetics, such as X-chromosome inactivation (XCI), the phenomenon where an entire X-chromosome is silenced in the cells of female mammals. Key questions and methods in the field are reviewed, including techniques to identify DNA sequences associated with different epigenetic modifications. Finally, we discuss how researchers are currently using these techniques to better understand the epigenetic regulation of gene function.

 Genetics

The Transwell Migration Assay

JoVE 5644

Cells migration in response to chemical cues is crucial to development, immunity and disease states such as cancer. To quantify cell migration, a simple assay was developed in 1961 by Dr. Stephen Boyden, which is now known as the transwell migration assay or Boyden chamber assay. This set-up consists an insert which separates the wells of a multiwell plate into top and bottom compartments. Cells whose migration is to be studied are seeded into the top compartment and the chemoattractant solution is placed in the bottom compartment. After incubation, counting the cells in the bottom compartment allows quantification of migration induced by chemoattractants.

This video will review the commonly used experimental set-up for cell migration studies. Then we'll highlight a few key considerations, and outline a generalized protocol for running an experiment involving adherent cells. Lastly, we'll review various adaptations of this set-up currently being used to study different factors that affect migration.

 Cell Biology

Two-Dimensional Gel Electrophoresis

JoVE 5686

Two-dimensional gel electrophoresis (2DGE) is a technique that can resolve thousands of biomolecules from a mixture. This technique involves two distinct separation methods that have been coupled together: isoelectric focusing (IEF) and sodium dodecyl sulfate polyacrylamide gel electrophoresis (SDS-PAGE). This physically separates compounds across two axes of a gel by their isoelectric points (an electrochemical property) and their molecular weights.

The procedure in this video covers the main concepts of 2DGE and a general procedure for characterizing the composition of a complex protein solution. Three examples of this technique are shown in the applications section, including biomarker detection for disease initiation and progress, monitoring treatment in patients, and the study of proteins following posttranslational modification (PTM).

Two-dimensional, or 2D, gel electrophoresis is a technique utilizing two distinct separation methods which can separate thousands of proteins from a single mixture. One of the techniques, SDS-PAGE or sodium dodecyl sulfate polyacrylamide gel electrophoresis, cannot fully separate complex mixtures alone. 2D gel electrophoresis couples the SDS-PAGE to a second method, isoelectric focusing or IEF, which separates based on isoelectric points, allowing for the resolution of potentially all proteins in a cell lysate. This video will show the principles of 2D gel electrophoresis, a general procedure, and some of its biomedical applications.

2D gel electrophoresis begins with IEF as the first dimension. Every protein has a pH value, called the isoelectric point or pI, where the net charge is zero. When a protein is subjected to an electric field, it will move toward the electrode with opposite charge. Samples of interest are loaded onto immobilized pH gradient, or IPG, strips which have embedded ampholytes, molecules containing both acidic and basic groups. An electric field is then applied to the pH gradient strip, causing the proteins to migrate until they reach the pH value matching their pI, where they lose their net charge.

Prior to running the second-dimension, the embedded proteins are treated with SDS, denaturing them and providing a uniform negative charge. Once completed, the IPG strips are placed onto a polyacrylamide gel. An applied electric field draws the proteins toward the anode, with larger proteins moving more slowly through the gel.

Once the protein mixture has been separated according to pI and molecular weight, the proteome map is visualized using stains, and proteins of interest are identified.

Now that we have discussed the principles of 2D gel electrophoresis, let's go over a typical laboratory procedure.

Before the experiment can be performed, the proteins must be solubilized into media. Solubilization of the sample is achieved by de-aggregating proteins with a combination of chaotropic agents for disrupting hydrogen-bonding interactions, nonionic detergents to prevent altering of the proteins' charge, reducing agents to break disulfide bonds, and buffers. To remove interfering abundant proteins and other molecules, the material is sequentially extracted by centrifugation, and collection of the resulting pellet; followed by treatment with endonuclease, an enzyme used to consume any DNA that would interfere with the experiment.

Once the proteins have been solubilized, IPG strips are prepared by rinsing with a strip-cleaning solution and left upside-down to dry. Each strip is then assigned a strip holder number. Once ready, the cellular extract is loaded onto the strips in a slow, sliding motion from the negative to the positive end. For the purpose of rehydration, damp blotting paper is placed on top of the electrode and under the gel strips; the IPG strips are then lined onto the IEF instrument. A high electric current is applied, and the proteins begin to migrate.

Following completion of the first dimension, the gel for SDS-PAGE is prepared in a casting apparatus. The IPG strips are treated by placing them face down in SDS-containing equilibration buffer. The electrophoresis unit is readied with the addition of electrophoresis buffer. The treated IPG strips are collected using tweezers, placed on top of the gel plates, and sealed with agarose-sealing solution. A voltage source then applies an electric field, which is held until the fastest-moving proteins are 1 cm from the bottom of the gel.

After completion of the electrophoresis, the proteins must be visualized. Traditionally this is performed by staining with Coomassie blue or silver nitrate. Proteins of interest may be transferred from the gel, and analyzed by Western blot analysis.

A second identification approach involves excision of the proteins from the gel, digesting them, than analyzing them by mass spectrometry.

Now that we've reviewed a procedure, let's look at some of the uses for 2D gel electrophoresis.

One of the most common uses for this technique is the identification of molecules involved in disease initiation and progression. 2D gel electrophoresis, coupled with mass spectrometry, can detect the up- or down-regulation of specific proteins in diseased areas in comparison to healthy ones.

Additionally, 2D gel electrophoresis is useful in following the progress of patients' response to a potential therapeutic drug. Specimens may be taken from patients at various timepoints following administration of the treatment. In this way, 2D gel electrophoresis coupled with Western blot or mass spectrometry analysis, can detect proteins associated with negative responses such as inflammation; or the absence of proteins in an alleviated state.

Another use for 2D gel electrophoresis is in the study of protein structure and function following posttranslational modification, or PTM, which are additions to proteins following their translation from mRNA. PTM's can regulate a variety of functions, including protein signaling, gene expression, or cause oxidative damage. 2D gel electrophoresis is sensitive to modifications such as methylation or acetylation, which can cause a shift in pI as well as the molecular weight.

You've just watched JoVE's video on 2D gel electrophoresis. This video described the principles of the technique, a typical experimental procedure, and several of its applications in the field of biomedicine.

Thanks for watching!

 Biochemistry

Collagen Hydrogels

JoVE 5786

Collagen is another widely used biomaterial that has found popularity in commercial applications, such as photography. Collagen has more recently been used in tissue engineering applications, by creating hydrogels that provide structure to engineered tissue.

This video introduces collagen as a biomaterial, demonstrates how it is harvested from porcine skin, and shows how the material is used to create a hydrogel for tissue engineering applications. Finally, several applications of the material and these techniques are shown.

 Bioengineering

Intermolecular Forces and Physical Properties

JoVE 11688

Intermolecular forces are attractive forces that exist between molecules. They dictate several bulk properties, such as melting points, boiling points, and solubilities (miscibilities) of substances. For example, a high-boiling-point liquid, like water (H2O, b.p. 100 °C), exhibits stronger intermolecular forces compared to a low-boiling-point liquid, like hexane (C6H14, b.p. 68.73 °C). The three kinds of intermolecular interactions include i) ion–dipole forces, ii) dipole–dipole interactions, and iii) van der Waals forces, which include London dispersion forces.

1. Ion–Dipole Forces

Ion–dipole forces are the electrostatic attractions between an ion and a dipole. They are common in solutions and play an important role in the dissolution of ionic compounds, like KCl, in water. The strength of ion–dipole interactions is directly proportional to i) the charge on the ion and ii) the magnitude of the dipole of polar molecules.

2. Dipole–Dipole Interactions

Polar molecules have a partial positive charge on one end and a partial negative charge on the other end of the molecule—a separation of charge called a dipole. The attractive force between two permanent dipoles is called a dipole–dipole attraction—the electrostatic force between the partially positive end of one polar molecule and the partially negative end of another. Hydrogen bonding is a type of dipole–dipole interaction between molecules with hydrogen, bonded to a highly electronegative atom, such as O, N, or F. The resulting partially positively charged H atom on one molecule (the hydrogen bond donor) could interact strongly with a lone pair of electrons of a partially negatively charged O, N, or F atom on adjacent molecules (the hydrogen bond acceptor). Hydrogen bonding increases the boiling point considerably.

3. van der Waals and London Dispersion Forces

The weakest of all forces is the van der Waals forces, which depend on the intermolecular distances between atoms and molecules. London dispersion forces, a subset of van der Waals forces, are experienced as a result of interactions between uncharged atoms/molecules owing to temporary, spontaneous shifts in electron distribution. The strength of these forces appears to increase with increasing molecular weight owing to the increase in surface area. As a result, compounds of higher molecular weights will generally boil at higher temperatures. Of note is that a branched hydrocarbon (neopentane) normally has a smaller surface area than its respective straight-chain (n-pentane) isomer, and therefore, a lower boiling point.

4. Solubility of Organic Compounds in Water

Liquids that can be homogeneously mixed in any proportion are said to be miscible. Miscible liquids have similar polarities. For example, methanol and water are both polar and capable of hydrogen bonding. On mixing, methanol and water interact through intermolecular hydrogen bonds of comparable strength to the methanol–methanol, and water–water interactions; thus, they are miscible. Likewise, nonpolar liquids like hexane and bromine are miscible with each other through dispersion forces. The chemical axiom “like dissolves like” is useful to predict the miscibility of compounds. Two liquids that do not mix to an appreciable extent are called immiscible. For example, nonpolar hexane is immiscible in polar water. Relatively weak attractive forces between the hexane and water do not adequately overcome the stronger hydrogen bonding forces between water molecules.

This text is adapted from Openstax, Chemistry 2e, Section 10.1: Intermolecular Forces, Section 11.3: Solubility, and Chapter 10: Liquids and Solids.

 Core: Organic Chemistry

Ethers from Alcohols: Alcohol Dehydration and Williamson Ether Synthesis

JoVE 11731

Overview

Ethers can be prepared from organic compounds by various methods. Some of them are discussed below,

Preparation of Ethers by Alcohol Dehydration

In this method, in the presence of protic acids, alcohol dehydrates to produce alkenes and ethers under different conditions. For example, in the presence of sulphuric acid, dehydration of ethanol at 413 K yields ethoxyethane, whereas it yields ethene at 443 K.

Figure1

This method is a nucleophilic substitution reaction. The two alcohol molecules involved in the reaction play two roles—one alcohol molecule acts as a substrate while the other acts as a nucleophile. The reaction follows an SN2 mechanism. The dehydration of secondary and tertiary alcohols to get corresponding ethers is unsuccessful as alkenes are formed easily in these reactions.

Preparation of Ethers by Williamson Ether Synthesis

It is the most versatile method for the preparation of asymmetrical ethers in laboratories. In this method, initially, the alcohol is deprotonated to form an alkoxide ion. Further, the alkoxide ion functions as a nucleophile and attacks an alkyl halide, leading to the formation of ether. The reaction generally follows the SN2 mechanism for primary alcohol.

Figure2

Williamson synthesis exhibits higher productivity when the halide to be displaced is on a methyl or a primary carbon. In the case of secondary alkyl halides, elimination competes with substitution, whereas the formation of elimination products is the only case in tertiary alkyl halides.

 Core: Organic Chemistry

Regioselectivity of Electrophilic Additions to Alkenes: Markovnikov's Rule

JoVE 11770

If a set of reactants can yield multiple constitutional isomers, but one of the isomers is obtained as the major product, the reaction is said to be regioselective. In such reactions, bond formation or breaking is favored at one reaction site over others.

The hydrohalogenation of an unsymmetrical alkene can yield two haloalkane products, depending on which vinylic carbon takes up the halogen. However, one product usually predominates, where hydrogen adds to the vinylic carbon bearing the greater number of hydrogens, and the negative part of the reagent adds to the more substituted vinylic carbon. Thus, one constitutional isomer is preferred, and the hydrohalogenation of alkenes is highly regioselective. Here, the regioselectivity is a consequence of the relative stabilities of the carbocation intermediates for each product.

The reaction begins with the transfer of a pair of electrons from the alkene to the proton of the hydrogen bromide, resulting in a carbocation. The protonation step is endergonic with a high energy of activation. It is also the slow rate-determining step. The rapid combination of the carbocation with the halide ion in the next step is exergonic with low activation energy.

Hammond's postulate indicates that the structure of the transition state in an endergonic process resembles that of the products as they are closer in energy. Thus, in the protonation step, the transition state structurally resembles the carbocation intermediate. Protonation can either yield the less substituted primary carbocation or, the more substituted tertiary carbocation. Tertiary carbocations are more stable than their secondary or primary counterparts, as their formation requires lower activation energy. Thus, the reaction prefers the pathway with the lower energy barrier via the more stable intermediate, resulting in regioselectivity.

When hydrohalogenation generates a new chiral center, the planar nature of the carbocation intermediate indicates that the nucleophile can approach from above or below the plane with equal probability, resulting in a racemic mixture of products.

When addition reactions proceed via mechanisms other than electrophilic addition, the hydrogen can add to the more substituted carbon, yielding the anti-Markovnikov product.

 Core: Organic Chemistry

Actin Filament Depolymerization

JoVE 11800

Actin filaments (F-actin) are composed of actin subunits. The dissociation of actin monomers can occur from either end of F-actin. The rate of dissociation is faster from the minus-end or the pointed end, where the actin subunits exist with a bound ADP, together known as ADP-actin. The depolymerization of F-actin is aided by proteins, including the actin-depolymerizing factor (ADF) and cofilin family of proteins, gelsolin, and glia maturation factor (GMF).

In F-actin, the ADF/cofilin proteins can bind with ADP-actins in a one-to-one ratio. The actin filament twists when these proteins bind to ADP-actin, generating mechanical stress and making the filament brittle. This stress allows rapid dissociation of the cofilin bound-ADP-G-actins from the filament. The ADF/cofilin proteins are also associated with AIP-1 (actin-interacting protein1), further enhancing the dissociation rate on the minus-end.

Gelsolin is a calcium ion-activated actin-binding protein. Activated gelsolin wedges into a straight actin filament and disrupts the interactions between the actin monomers within the filament, splitting the F-actin into two. One end has a newly formed gelsolin-capped plus-end, while the other filament has a newly formed minus-end with rapidly dissociating ADP-G-actins.

The disassembly of actin filaments can also occur at branched filaments bound to the Arp2/3 complex. GMF binds to the Arp2/3 complex at the branch junction of actin networks. Upon binding, it prevents further nucleation of actin filaments from that site. GMF plays an important role in the lamellipodia formation required for cell movement and migration.

 Core: Cell Biology

Microtubule Associated Proteins (MAPs)

JoVE 11909

Microtubule function and architecture are regulated by an array of specialized proteins called microtubule-associated proteins or MAPs. These proteins are widespread across different organisms and have conserved protein motifs, like the multi-TOG domain for tubulin binding found in the CLASP family of MAPs. Some MAPs are lineage-specific based on their conserved domains. Their functions depend upon the cytoskeletal architecture and cell type they are located within. In-plant cells, a specific microtubule-associated protein‒ tortifolia, binds with cortical microtubules to regulate organ the orientation and direction of organ growth. On the other hand, tau proteins are specifically associated with microtubules in neurons in animal cells. MAPs were first identified within neurons and were named “classical MAPs''.

Depending on how MAPs regulate microtubules, they are broadly classified as stabilizers, destabilizers, capping proteins, crosslinkers, and cytoskeleton integrator proteins.  MAPs are further classified based on where they localize on the microtubules. They are broadly divided into three groups: Lattice-binding proteins, microtubule plus-end trafficking proteins, and minus-end targeting proteins. Lattice binding proteins bind along the filament length instead of the microtubule plus or minus end. Tau and MAP2 found in neurons' axonal and dendritic microtubules belong to the lattice-binding MAPs. Microtubule plus-end trafficking proteins include proteins that target the growing end of the microtubules. Examples include EB1, XMAP-215, and kinesin-13. EB1 and XMAP-215 are microtubule-stabilizing and growth-promoting proteins, while kinesin-13 is a microtubule destabilizer. Minus-end targeting proteins include microtubule formation initiator proteins like the γ-tubulin ring complex (γ-TRC) and capping proteins like the calmodulin regulated spectrin-associated protein family (CAMSAPs) members. CAMSAPs bind to the minus-end of microtubules to stabilize them and prevent the dissociation of the tubulin subunits.

The dynamic structure of microtubules varies throughout the cell cycle. During interphase, the microtubule network transports organelles and vesicles and helps organize the cytoskeleton within the cell. As the cell enters into a dividing phase, the previous microtubule mesh disassembles and reorganizes into mitotic spindles that aid in separating chromosomes and cytokinesis. These functions and variability of microtubules are possible due to the various MAPs present within the cell.

 Core: Cell Biology

Conversion of Alcohols to Alkyl Halides

JoVE 11929

This lesson delves into the conversion of alcohols to corresponding alkyl halides and the mechanism of action for different reagents. Typically, the hydroxyl group is first protonated to convert it to a stable leaving group. Consequently, based on the starting alcohol, the mechanism undergoes either of the nucleophilic substitution routes, SN1 or SN2. Tertiary alkyl halides are made using the two-step SN1 mechanism that occurs via a carbocation intermediate, which is stabilized by hyperconjugation. However, for primary alcohols, the protonation of the hydroxyl group leads to the concerted SN2 route. Secondary alcohols can proceed via either mechanism based on the reaction conditions.

The popular reagents used for converting alcohols to corresponding alkyl halides include the hydrogen halides like hydrogen bromide and hydrogen chloride. However, while it is straightforward with the former, the latter needs an additional catalyst like zinc chloride. This catalyzes the hydroxyl group into a better leaving group enabling the subsequent SN2 process. Other reagents of choice are thionyl chloride and phosphorus tribromide with a similar mechanism. In the presence of relatively weak bases like pyridine/tertiary amine, they generate an excellent leaving group compared to the original leaving group of water.

The most exciting class of reagents is sulfonyls. They react with the alcohols to form corresponding mesylates, tosylates, or triflates to improve their reactivity in an SN2 reaction. In these species, resonance stabilization is inherent to the sulfonyl group. Additional resonance stabilization is contributed by the benzene ring of the tosyl group, and further stability is provided by the strongly electron-withdrawing trifluoromethyl in the triflate.

Stereochemistry

Most importantly, the choice of reagent influences the stereochemistry of the product formed. The use of thionyl chloride leads to an inversion of configuration, while tosyl chlorides retain the chiral configuration in the native alcohol.

 Core: Organic Chemistry

Cancer Cell Migration through Invadopodia

JoVE 12254

Invadosome is a broad category of cell surface structures with proteolytic activity that  degrades the extracellular matrix (ECM). Invadosomes are present in normal cell types, including macrophages, endothelial cells, and neurons, as well as tumor cells. Although the macrophage podosomes and tumor cell invadopodia are classified as invadosomes, they have different structures, molecular pathways, and functions. Podosomes are short structures that last for a few minutes. However, invadopodia can last for hours and are many microns in length.

In tumors, invadopodia are essential for cell intravasation and extravasation through blood vessels. They are similar to lamellipodia and filopodia, with a filamentous actin core shaped by actin nucleators and regulatory factors. The lifecycle of an invadopodium is initiated by complex processes involving various signaling cascades that remodel the cytoskeleton and cell membrane. The actin cytoskeleton is reorganized and assembled into new filaments and branches, which requires the activity of actin nucleators such as formins or the ARP2/3 complex.

The Arp2/3 complex requires nucleation-promoting factors (NPF) for optimal activity as it is slow to spontaneously initiate new actin branches. Cortactin, an NPF and scaffold protein, recruits ARP2/3 to filaments, allowing branching of the actin network. Additionally, cortactin can stabilize the newly generated branches. The formin proteins induce the elongation of unbranched actin filaments, while filament bundling is coordinated by fascin. Together these events promote actin filament polymerization and maturation of the invadopodium. 

After the invadopodium stabilizes, kinesins use the neighboring microtubules as tracks to transport vesicles with proteases from the Golgi network. These vesicles release enzymes such as  matrix metalloproteases, cathepsins, and serine proteases from the membrane of the invadopodia to degrade the surrounding ECM. Lastly, the actin core is disassembled to retract the invadopodium. Several proteins are implicated in the retraction pathway, but it is unclear how these proteins interact during the final stage of the invadopodium life cycle. 

 Core: Cell Biology

orthopara-Directing Activators: –CH3, –OH, –⁠NH2, –OCH3

JoVE 12472

All ortho–para directors, excluding halogens, are activating groups. These groups donate electrons to the ring, making the ring carbons electron-rich. Consequently, the reactivity of the aromatic ring towards electrophilic substitution increases. For instance, the nitration of anisole is about 10,000 times faster than the nitration of benzene. The electron-donating effect of the methoxy group in anisole activates the ortho and para positions on the ring and stabilizes the corresponding intermediates through resonance effects via pi-donation. As a result, the energy of the transition state is lowered for ortho and para intermediates, leading to an accelerated reaction.

 Core: Organic Chemistry

Transcytosis of IgG

JoVE 12554

Transcytosis is the process in which molecules are internalized by endocytosis, transported across the cell, and released through exocytosis from the opposite end of the cell. Molecules such as insulin, immunoglobulins, and certain nutrients are transferred through the recycling endosomes by recycling and transcytosis.

IgG molecules from a mother undergo transcytosis starting around 13 weeks of gestation. The amount of IgG transferred and entering the fetal blood circulation increases with gestational age. In the process, IgG molecules cross two cellular barriers—syncytiotrophoblasts of the placental villi and the endothelial cells of the fetal capillaries.

The placental transfer of IgG from the mother to the developing fetus offers a protective advantage to both mother and fetus from infections. Specific vaccines taken during pregnancy, e.g. tetanus toxoid or inactivated influenza vaccines, elicit the production of IgG antibodies, which, after transcytosis, can potentially protect the neonate.

IgA undergoes transcytosis through the intestinal epithelial cells and is released into the intestinal lumen. IgA is produced by plasma cells as a monomer and later dimerizes. The dimeric IgA binds to the poly-IgA receptor (pIgR) at the basolateral side of the intestinal epithelial cells. The IgA-receptor complex first enters the intestinal epithelial cells through endocytosis, fuses with a recycling endosome, and is released by exocytosis at the apical end of the intestinal epithelial cells. This trajectory of IgA is opposite to that followed by IgG. In the transfer of IgG, endocytosis occurs at the apical end and exocytosis occurs at the basolateral end.

 Core: Cell Biology

Introduction to Fibroblasts

JoVE 12823

Rudolph Virchow discovered spindle-shaped cells called fibroblasts in 1858. Inactive fibroblasts, called fibrocytes, become activated by various stimuli, such as growth factors and inflammatory cytokines. Activated fibroblasts play a crucial role in wound healing, inflammation, formation of new blood vessels, and cancer progression. Uncontrolled activation of fibroblasts results in fibrosis, the excess deposition of fibrous tissue, which can lead to scarring and affect normal organs. This results in fibrotic disorders such as liver cirrhosis, kidney cirrhosis, and cardiac fibrosis.

In addition to being easily accessible in the body, fibroblasts can also be cultured in the laboratory as primary cell cultures or permanent cell lines. Fibroblast cell lines have been used for years to determine the pathogenesis of certain specific diseases. Currently, fibroblasts are also used for modeling diseases.

Fibroblasts can retain the memory of their anatomical positions through changes in gene expression and chromatin modifications. This includes memories of the location of their tissue of origin and memories of previous inflammation.  This property of keeping memories of various stimuli supports fibroblasts’ role in immune responses and tissue homeostasis.

 Core: Cell Biology

IP3/DAG Signaling Pathway

JoVE 13326

Membrane lipids such as phosphatidylinositol (PI) are precursors for several membrane-bound and soluble second messengers. Specific kinases phosphorylate PI and produce phosphorylated inositol phospholipids. One such inositol phospholipids are the  phosphatidylinositol-4,5 bisphosphate [PI(4,5)P2], present in the inner half of the lipid bilayer. Upon ligand binding, GPCR stimulates Gq proteins to turn on phospholipase Cꞵ. Activated phospholipase Cꞵ cleaves PI(4,5)P2 and produces two-second messengers—a membrane-bound diacylglycerol (DAG) and a cytosolic inositol-1,4,5 trisphosphate (IP3). The events mediated by these second messengers are called IP3/DAG pathway.

IP3 is a sugar-phosphate molecule. They are soluble and can rapidly diffuse through the cytosol to reach the endoplasmic reticulum (ER). IP3 binds and opens IP3-gated calcium channels on the ER membrane and drives out calcium into the cytosol. This way, IP3 transmits an external signal by increasing the cytosolic concentration of another second messenger, like the calcium ions. IP3 triggers various cellular responses via elevating cytosolic calcium levels, including smooth muscle contraction in the blood vessels and platelet aggregation. However, IP3 also controls the rising cytosolic calcium levels. IP3 is rapidly degraded to inositol-1,4 bisphosphate, which cannot bind or open ER calcium channels. This prevents calcium release to the cytosol.

The second product, the membrane-bound second messenger DAG also plays an essential role in various cellular processes. They bind and activate protein kinase C (PKC), which is involved in cellular growth and metabolism. PKC phosphorylates multiple transcription factors, which move to the nucleus and initiate the transcription of genes involved in cell division.

DAG can also be cleaved to form arachidonic acid,  a precursor for eicosanoids, a small lipid signal molecule. A commonly known eicosanoid, prostaglandin, affects pain and inflammatory responses. Many anti-inflammatory and pain-relieving drugs available commercially, like aspirin, ibuprofen, and cortisone, inhibit prostaglandins synthesis.

 Core: Cell Biology

Cell Adhesion in Plants

JoVE 13370

Plants have rigid cell walls that are made up of cell wall polysaccharides that mediate cell-cell adhesion. The primary cell walls of plants consist of two independent and interacting polysaccharide networks: a pectin matrix that embeds the second network comprising cellulose and hemicelluloses.

Pectins are complex heteropolymers mainly composed of negatively-charged α-D-glucopyranosyl uronic acid and some neutral glycosyl residues such as α-L-rhamnopyranose, α-L-arabinofuranose, and β-D-galactopyranose. Pectin networks crosslink one cell’s middle lamella to others, thus acting as the most important tool for cell-cell adhesion.

Adjacent cell walls are connected via membrane-lined channels called plasmodesmata (singular plasmodesma). Each plasmodesma allows the transport of molecules via the connected cytoplasm through contact-dependent signaling, like gap junctions in animals. Unlike gap junctions, plasmodesmata are more flexible because they let molecules pass through the cell wall and membrane. The plasmodesmata also allow direct communication from a single cell to many cells, joining the cells in a symplast network.

Other cellular components such as ferulic acids, xyloglucan-like polysaccharides, and proteins such as specialized wall-associated kinases and extensins also regulate cell adhesion during growth and development, although the exact mechanisms remain unknown.

 Core: Cell Biology

Applications Of NMR In Biology

JoVE 13387

Nuclear magnetic resonance (NMR) spectroscopy is a very valuable analytical technique for researchers. It has been used for more than 50 years as an analytical tool. F. Bloch and E. Purcell formulated NMR in 1946 and won the 1952 Nobel Prize in Physics  for their work. Biological macromolecules such as proteins, nucleic acids, lipids, and organic molecules including pharmaceutical compounds, can be studied using this versatile tool that exploits the magnetic properties of certain nuclei.

The basic principle of this technique is that nuclei, in addition to their electric charge, also have a spin. Nuclei with an odd atomic number or mass possess this property of spin, which is vital for the NMR technique. The spin is random and in random directions,  similar to a spinning top. Hence, when placed under an external magnetic field, these nuclei align themselves with or against the applied field. These nuclei return to their original orientation when the external field is removed.  The energy gap is then translated into a spectra that depends on the nature of the environment of the atoms, and the distance between nuclei. The resulting spectra helps study different parameters like the structure, dynamics, and properties of the samples. Properties such as  reaction state,chemical environment, and interactions of the samples are examples of investigational results that can be studied with this technique.

In biology, 13C, 1H, 2H 15N, 31P, 23Na, and 19F are important biologically relevant NMR-active nuclei that help understand biochemical pathways involved in amino acid, lipid, and carbohydrate metabolism. Also, NMR offers a window into observing and quantifying numerous compounds in biological fluids, cell extracts, and tissues without the need for complex sample preparation or fractionation.

Over the past two decades, NMR has been developed to produce detailed images in a process now called magnetic resonance imaging (MRI), a name coined to avoid the use of the word “nuclear” and the concomitant implication that nuclear radiation is involved. MRI is based on NMR, in which an externally applied magnetic field interacts with the nuclei of certain atoms, particularly those of hydrogen (protons) of the body tissue.

Numerous applications of this technique, including its pivotal role in drug discovery and proteomics, are helping  advance research to new heights, benefiting humanity.

 Core: Cell Biology

Immunogold Electron Microscopy

JoVE 13403

Immunoelectron microscopy utilizes immunogold labeling of endogenous proteins with specific antibodies to detect and localize these proteins in cells and tissues. The procedure provides insights into the distribution and quantification of protein under different stimulation conditions offering clues about their functions. Conjugating highly electron-dense gold particles with primary or secondary antibodies allow antigen detection on and within cells, with high resolution and specificity. Immunoelectron microscopy has been used to identify the cellular and subcellular localization of proteins involved in neurotransmission, nuclear protein components, and identification of immune cell types.

Three approaches are applied to localize cell antigens using transmission electron microscopy. First, when studying the localization of intracellular antigens, protocols involve antibody labeling post-embedding in acrylic resins, antibody labeling pre-embedding in the resin combined with cell membrane permeabilization, or cryo-ultramicrotomy without embedding. Second, when cell-surface proteins are to be localized, the pre-embedding protocol is used where labeling with immunogold antibodies is done before embedding in the resin. The third technique, the whole-mount technique, does not involve resin-embedding, and the antibody reactions are carried out directly on a sample to locate the surface molecules. 

For scanning electron microscopy, the surface of interest must necessarily be exposed. The sample can be directly fixed and labeled if the outer surface of the plasma membrane is of interest. However, the cells must be permeabilized with detergents and immunolabeled for studying intracellular components or antigens integrated into a tissue. For viewing, the sample is freeze-fractured (rapid freezing and breaking with a knife). Many intracellular components, particularly macromolecular complexes, like chromatin, protein complexes, or viruses, can be isolated, fixed, and immunolabeled for SEM.  Finally, an immunonegative staining technique is applied for surface antigens on small specimens, such as viruses and bacteria, which lend themselves to negative staining.

 Core: Cell Biology

Multipotency and Niche of Bulge Stem Cell

JoVE 13467

A hair follicle or HF is a small part of the skin that produces the hair shaft. Paul Gerson Unna was the first to observe a bulge in the human hair follicle's outer root sheath (ORS). The bulge is present between the sebaceous gland and the arrector pili muscle and is the niche for hair follicle stem cells (HFSCs). The bulge is also a niche for melanocyte stem cells, and their loss results in graying of hair. The HFSCs express Sox9 and Lhx2, which help them maintain stemness and prevent differentiation. If Sox9 is inhibited, the cells begin to differentiate.

The HF undergoes three cyclic phases: growth phase or anagen, transitional phase or catagen, and resting phase or telogen. During the anagen phase, HFSCs and hair germ cells are activated. Hair germ cells proliferate, form the ORS, and move upwards, differentiating into follicle cells, and the hair grows. In the catagen phase, hair germ cells undergo apoptosis, hair stops growing, and only bulge stem cells survive. Telogen is the resting phase where HFSCs undergo quiescence and the hair rests in the follicle. 

Adult stem cells are identified using label-retaining techniques where the DNA of cells is labeled with bromodeoxyuridine (BrdU) or tritiated-thymidine (3 H-T). The cells which divide very slowly contain the labels, while fast proliferating cells do not retain the labels. The labeled, slow-cycling cells are called label-retaining cells (LRCs). Bulge stem cells are LRCs.

 Core: Cell Biology

Forced Transdifferentiation

JoVE 13484

Transdifferentiation, also known as lineage reprogramming, was first discovered by Selman and Kafatos in 1974 in silkmoths. They observed that the moths’ cuticle-producing cells transformed into salt-producing cells. Many such cases of natural transdifferentiation occur in organisms. In humans, pancreatic alpha cells can become beta cells. In newts, the loss of the eye’s lens causes the pigmented epithelial cells to transdifferentiate into the lens cells.

Artificial transdifferentiation occurs when a transcription factor is forced to be expressed in a mature cell type. Therapeutically, turning on the transcription factor, MyoD converts human fibroblasts into muscle cells. Another transcription factor, C/EBP α, transforms mature lymphocytes into macrophages. In rats, the chemical dexamethasone turns on the expression of a transcription factor, C/EBP β, to convert pancreatic exocrine cells into liver cells. This transdifferentiation likely occurs since both pancreatic and liver cells originate from neighboring regions of the endoderm during development. The reverse of this differentiation, i.e., conversion of liver cells into pancreatic exocrine cells, is also experimentally possible. 

Transdifferentiation of cells does not go through an intermediate pluripotent stem cell state; therefore, it has some advantages over differentiating cells from induced pluripotent stem cells (iPSCs). Some epigenetic marks must be removed while transdifferentiating cells, while all the epigenetic marks are erased while reprogramming iPSCs. Mutations are less likely to occur when cells directly transform from one mature type to another.

 Core: Cell Biology

Detection of Gross Error: The Q Test

JoVE 14517

When one or more data points appear far from the rest of the data, there is a need to determine whether they are outliers and whether they should be eliminated from the data set to ensure an accurate representation of the measured value. In many cases, outliers arise from gross errors (or human errors) and do not accurately reflect the underlying phenomenon. In some cases, however, these apparent outliers reflect true phenomenological differences. In these cases, we can use statistical methods to help determine whether to retain the outlier. A statistical method that can help us retain or reject outliers is called the Q-test. To perform the Q-test, we first arrange the values in a data set in order of increasing value. Then, we calculate the Q value by taking the ratio of the absolute difference between a data point and its adjacent data point. This Q value is then compared with the tabulated critical Q value at a chosen significance level and appropriate degrees of freedom. If the Q value equals or exceeds the reference Q value in the table, the data point is considered an outlier and therefore rejected from the data set. Here, it is reasonable to disregard the data point as an outlier because the magnitude of the deviation cannot be logically accounted for by random (indeterminate) errors. On the other hand, if the Q value is smaller than the reference value in the table, the data point should be retained, and the interpretation is that the difference between this data point and the rest of the data is within reasonable (statistical) expectation.

 Core: Analytical Chemistry

Ladder Diagrams: Complexation Equilibria

JoVE 14533

Ladder diagrams are useful for evaluating equilibria involving metal-ligand complexes. The vertical scale of the ladder diagram represents the concentration of unreacted or free ligand, pL. The horizontal lines on the scale depict the log of stepwise formation constants for metal-ligand complexes and indicate the dominant species in all the regions.

The formation constant, K1, for the formation of Cd(NH3)2+ complex from cadmium and ammonia is 3.55 × 102. Log K1 (i.e. pNH3) is 2.55, and represents the dividing line between the predominance regions for Cd2+ and Cd(NH3)2+. Above the value of 2.55, Cd2+ is the predominant species.

Alternatively, ladder diagrams of complexation reactions can also be constructed using cumulative formation constants instead of stepwise formation constants. For example, the ladder diagram for the Zn2+-NH3 system uses the cumulative formation constants, showing [Zn(NH3)4]2+ as the dominant species at lower pNH3 values. At higher pNH3 values, Zn2+ predominates.

 Core: Analytical Chemistry

EDTA: Chemistry and Properties

JoVE 14572

Polydentate ligands are most widely used in complexometric titrations because they form more stable complexes with the metal ions than mono- or bidentate ligands due to the chelate effect. Examples of polydentate ligands are ethylenediaminetetraacetic acid (EDTA), crown ethers, and cryptands. The most important feature of optimal polydentate ligands is the ability to form 1:1 complexes in a single-step process. Amino carboxylic acid derivatives are frequently used as complexing agents. EDTA is among the most versatile ligands, and it is widely used as a chelating agent in analytical chemistry because most elements can be measured with EDTA through different titration methods.

EDTA is a hexadentate ligand consisting of six complexing groups: four carboxylate oxygens and two amine nitrogens. They coordinate with the metal ion by sharing the lone pairs with the metal. EDTA forms stable cage-like structures with most metal ions and is usually represented with octahedral geometry. However, the geometry of the different complexes depends on the size of the metal ions.

EDTA is a neutral tetraprotic acid, represented by the abbreviation H4Y. (The disodium salt of EDTA is the preferred laboratory reagent due to its higher solubility than the parent acid.) EDTA dissociates into various species, and their relative amounts depend on the pH of the solution. At low pH, the amine nitrogens are protonated, making EDTA a hexaprotic system (H6Y2+). The fully deprotonated form (Y4−) exists at high pH and can generate hexadentate complexes with metal ions. This is considered the 'active' form of EDTA.

 Core: Analytical Chemistry

Washing, Drying, and Ignition of Precipitates

JoVE 14588

After filtration, the precipitate is washed to remove coprecipitated impurities and any remaining mother liquor. Colloidal precipitates, such as silver chloride, are washed with an electrolyte (such as dilute nitric acid) to prevent the peptization of the precipitate. In the case of slightly soluble precipitates, the wash solution contains a common ion to reduce solubility. Lead sulfate, which is slightly soluble in water, is washed with dilute sulfuric acid. Similarly, wash solutions may be basic or acidic to prevent the hydrolysis of salts of weak acids and weak bases, respectively. For example, magnesium ammonium phosphate hexahydrate is washed with a dilute ammonia solution. Following washing, the precipitate is dried to remove adsorbed water and electrolyte. In some cases, ignition is required to convert the precipitate to a suitable weighing form. For example, heating at 900 °C converts magnesium ammonium phosphate  hexahydrate to magnesium pyrophosphate for weighing. Here, the excess water that may or may not be in the crystals depending on the levels of heating or drying prevents stoichiometric analysis. Finally, the precipitate is cooled in a desiccator and weighed.

 Core: Analytical Chemistry

Extraction: Effects of pH

JoVE 14621

Consider a neutral form of an amine, B, with a partition coefficient, K, in a liquid mixture containing organic and aqueous phases. The pH of the aqueous phase affects the charge on acidic and basic solutes, and the charged form is usually more soluble in the aqueous phase. Suppose the conjugate acid form of the amine is soluble only in the aqueous phase while the base form is soluble in both phases. Then the distribution coefficient, D, can be given as the ratio of amine concentration in the organic phase to the sum of concentrations of the amine and its conjugated acid in the aqueous phase. In this equation, the terms for partition coefficient and the acid dissociation constant, Ka, can be substituted with the concentrations, and the fraction of neutral amine in the aqueous phase can be determined by solving the equation. With ladder diagrams and a plot of log D versus pH, it can be concluded that a neutral base can be extracted into the aqueous phase at a pH low enough to convert it into its charged conjugate acid. Similarly, a charged acid can be extracted into the organic phase at a pH high enough to convert it to its neutral conjugate base.

 Core: Analytical Chemistry

Types of Skeletal Muscle Fibers

JoVE 14850

Skeletal muscles comprise various fibers, each with distinct characteristics and roles in movement and stability. They are mainly categorized into three types — fast-twitch, slow-twitch, and intermediate.

Fast-twitch fibers

Fast-twitch fibers, or Type II fibers, are designed for quick, powerful bursts of speed and strength. They reach peak tension within approximately 0.01 seconds following stimulation. Characterized by a large diameter and densely packed myofibrils, these fibers contain abundant glycogen reserves for rapid energy access. However, their relatively low mitochondrial content limits their endurance, making them prone to quick fatigue. Fast-twitch fibers primarily rely on anaerobic metabolism, consuming high levels of ATP quickly, and are best suited for short, intense activities like sprinting or heavy lifting.

Slow-twitch fibers

In contrast, slow-twitch fibers, or Type I fibers, are the endurance powerhouses of the muscle. With a smaller diameter, these fibers are packed with capillaries, high levels of myoglobin, and numerous mitochondria. This combination allows for efficient oxygen transport and utilization, supporting prolonged aerobic activity and ATP production. Slow-twitch fibers contract more slowly and can sustain contractions for extended periods, making them ideal for endurance activities like long-distance running or cycling. They are less reliant on anaerobic metabolism due to their enhanced oxygen reserves and blood supply.

Intermediate fibers

Intermediate fibers bridge the gap between fast and slow-twitch fibers. They resemble fast-twitch fibers in appearance, with minimal myoglobin content and a lighter color. However, they boast a more developed capillary network and a higher mitochondrial density than fast-twitch fibers. This endows them with greater resistance to fatigue and a capacity for both anaerobic and aerobic energy production. Intermediate fibers are adaptable and can change their characteristics with endurance or strength training, making them versatile components in various physical activities.

 Core: Anatomy and Physiology

Muscles of the Vertebral Column

JoVE 14871

The back muscles that lie deep into the thoracolumbar fascia are called intrinsic or true back muscles. These muscles are divided into four layers: superficial, intermediate, deep, and deepest layers.

Superficial Layer:

The superficial layer consists primarily of the splenius muscles, which include the splenius capitis and splenius cervicis. These muscles are mainly responsible for the head and cervical spine movements, including extension, rotation, and lateral bending. The splenius capitis extends from the upper thoracic and lower cervical vertebrae to the skull, while the splenius cervicis stretches between the thoracic and cervical vertebrae.

Intermediate Layer:

The intermediate layer includes the erector spinae muscles, a group of three long, column-like muscles: the iliocostalis, longissimus, and spinalis. These muscles run almost the entire length of the spine, from the sacrum and iliac crest to the ribs and vertebrae up to the base of the skull. They are crucial for maintaining posture and providing the strength to extend the vertebral column. The erector spinae muscles also assist in lateral flexion and, when acting unilaterally, aid in bending the spine to the side.

Deep and Deepest Layers:

The deep layer comprises the transversospinalis muscles, including the semispinalis, multifidus, and rotatores. These shorter muscles are primarily involved in stabilizing the spine, aiding in extension, and allowing for controlled rotation and bending of the vertebral column. The deepest layer consists of minor deep muscles like the interspinales and intertransversarii, which span between adjacent vertebrae and contribute to the stability and localized movements of the spine. Additionally, the levatores costarum, found in this layer, assist in elevating the ribs, which is particularly important during deep inhalation.

 Core: Anatomy and Physiology

Nervous Tissue: Neuron Types

JoVE 14887

Neurons, the fundamental units of the nervous system, can be classified based on both their structural and functional characteristics.

Structurally, neurons are categorized into three main types: multipolar, bipolar, and unipolar (or pseudounipolar). Multipolar neurons, which are the most common type in the brain and spinal cord, as well as all motor neurons, possess multiple dendrites and a single axon.

Bipolar neurons, on the other hand, have one primary dendrite and one axon. They are typically located in specialized sensory areas such as the retina of the eye, the inner ear, and the olfactory region of the brain. These neurons play crucial roles in sensory perception, processing information about vision, hearing, and smell.

Unipolar or pseudounipolar neurons are unique in that their dendrites and axons fuse together into a single process extending from the cell body. This structure results from neuronal development, where the neuron begins as a bipolar neuron and later merges its dendrites and axon. The dendrites of unipolar neurons often function as sensory receptors, detecting stimuli like touch, pressure, pain, or temperature changes.

Functionally, neurons can be divided into sensory neurons, motor neurons, and interneurons. Sensory neurons, which are primarily unipolar and occasionally bipolar, carry information from sensory receptors to the central nervous system (CNS). Multipolar motor neurons convey commands from the CNS to peripheral effectors such as muscles. Interneurons, also multipolar, serve as intermediaries, transmitting information between sensory and motor neurons.

 Core: Anatomy and Physiology

Anatomy of the Brain: Ventricles

JoVE 14904

There are hollow fluid-filled cavities known as ventricles deep inside the human brain. There are two lateral ventricles, one in each cerebral hemisphere, and each has three different projections — the anterior, inferior, and posterior horns visible from the lateral side. A thin membrane called the septum pellucidum separates the two lateral ventricles. The slender third ventricle in the diencephalon is connected to each lateral ventricle via a channel called the interventricular foramen. The cerebral aqueduct is posterior to the third ventricle, a canal-like structure that connects it to the fourth ventricle between the pons and the anterior surface of the cerebellum.

The brain ventricles are filled with cerebrospinal fluid (CSF), produced and filtered from a network of blood capillaries called the choroid plexus. These capillaries are lined by tightly joined ependymal cells that secrete CSF into the ventricles. The CSF then drains into the central canal of the spinal cord and surrounds the brain. The CSF mainly provides mechanical protection to the central nervous system, enables optimal neural signaling, and acts as a medium for nutrient exchange.

 Core: Anatomy and Physiology

Peripheral Nervous System: Ganglia and Nerves

JoVE 14920

The Peripheral Nervous System (PNS) is a crucial component of the body's neural network, extending beyond the central nervous system (CNS) to bridge the gap between the CNS and the external environment. It encompasses nerves, ganglia, and sensory receptors.

Nerves

The nerve is a bundle of axons that serves as the communication highway in the PNS. Each nerve is ensheathed in a protective layer of connective tissue called the epineurium. This outermost layer safeguards the nerve and supports the blood vessels supplying it. Within the nerve, axons are organized into smaller bundles called fascicles, each surrounded by the perineurium. This arrangement facilitates the efficient transmission of electrical signals along the axons. Encasing each axon and its Schwann cells is a delicate sheath called the endoneurium. This layered structure ensures a precise and protected transmission of nerve impulses across the PNS.

The PNS is functionally divided into two principal categories — the cranial and spinal nerves. The 12 pairs of cranial nerves have their origins in the brain. They are instrumental in various functions ranging from sensory input — such as vision, smell, and taste — to motor control, which includes facial expressions, eye movement, and speech.

On the other hand, the 31 pairs of spinal nerves emerge from the spinal cord, branching out to innervate the entire body. These nerves transmit sensory information from the limbs and trunk to the CNS and carry motor commands from the CNS to the muscles. Each spinal nerve is a mixed nerve, containing afferent (sensory) and efferent (motor) fibers, allowing for the bidirectional information flow between the CNS and the body.

Ganglia

Ganglia are clusters of neuron cell bodies located outside the CNS that are integral to the PNS's operation. They are categorized based on their location and the specific functions they perform.

The dorsal root ganglia, situated along the spinal cord, house the cell bodies of sensory neurons. They relay sensory information from peripheral receptors to the spinal cord, facilitating the sense of touch, pain, temperature, and proprioception.

The autonomic ganglia are associated with the motor division of the PNS. They play a key role in regulating involuntary functions, such as digestion, respiration, and heart rate. These ganglia are subdivided into sympathetic and parasympathetic ganglia, reflecting their roles in the body's fight-or-flight response and rest-and-digest activities, respectively. Sympathetic ganglia are positioned near the spinal cord, ready to mobilize the body's resources under stress. In contrast, parasympathetic ganglia are located close to visceral organs, overseeing the conservation and restoration of the body's energy.

 Core: Anatomy and Physiology

Indirect Motor Pathways

JoVE 14938

The indirect motor or extrapyramidal pathways originate in the brainstem, the lower portion of the brain that connects it to the spinal cord. They consist of several distinct tracts, each with specialized functions. The four main tracts of the indirect motor pathways are the vestibulospinal tract, the reticulospinal tract, the tectospinal tract, and the rubrospinal tract.

The vestibulospinal tract originates in the vestibular nuclei of the brainstem. The vestibular system detects changes in head position and orientation, which is critical for maintaining balance and posture. The vestibulospinal tract receives input from the vestibular system about the position of the head and from the proprioceptors and eyes about the position of muscles, joints, and body in space. By regulating muscle tone, the vestibulospinal tract helps to maintain an upright posture, allowing the body to remain stable in a gravitational field.

The reticulospinal tract also originates in the brainstem, specifically in the reticular formation. This tract has two distinct components: the lateral reticulospinal and medial reticulospinal tracts. The lateral reticulospinal tract is responsible for inhibiting the muscles of the trunk and proximal limbs, while the medial reticulospinal tract excites these same muscles. By balancing the excitatory and inhibitory signals, the reticulospinal tract helps to maintain muscle tone and balance during ongoing movements, such as lifting weights.

The tectospinal tract originates in the superior colliculus, a region of the brainstem involved in processing visual and auditory information. This tract controls sudden head and neck movements in response to visual and auditory stimuli. Additionally, it plays a role in controlling rapid eye movements, which are critical for clear vision during rapid head movements.

Finally, the rubrospinal tract originates in the red nucleus. This tract controls precise and voluntary movements of the distal part of the upper limbs, including the hands and fingers. This tract is essential for tasks that require fine motor control, such as playing a musical instrument or typing on a keyboard.

In summary, the indirect motor pathways are critical for the control of movement throughout the body. By regulating muscle tone, maintaining balance and posture, and allowing us to respond to visual and auditory stimuli, these pathways play a critical role in allowing a person to move efficiently and precisely.

 Core: Anatomy and Physiology

Parasympathetic Signaling

JoVE 14955

Parasympathetic signaling plays a crucial role in regulating various physiological processes. It involves the release of acetylcholine (ACh) by parasympathetic neurons, which can have localized and short-lived effects. The majority of ACh released is rapidly inactivated at the synapse by the enzyme acetylcholinesterase (AChE), which hydrolyzes Ach into choline and acetate. Additionally, the tissue cholinesterase deactivates any ACh diffusing into the surrounding tissues.

The effects of parasympathetic stimulation are tightly controlled and depend on the specific receptors involved, allowing for precise modulation of target organs and tissues. There are two main types of cholinergic receptors: nicotinic and muscarinic receptors.

Nicotinic Receptors

Nicotinic receptors are found on postganglionic cells of both the parasympathetic and sympathetic divisions and at neuromuscular junctions in the somatic nervous system. Activation of these receptors by ACh leads to the opening of chemically gated sodium (Na+) channels in the postsynaptic membrane of the postganglionic neuron or motor end plate of the muscle fiber. This excitation results in the activation of the respective cell.

Muscarinic Receptors

Muscarinic receptors are found at cholinergic neuromuscular or neuroglandular junctions in the parasympathetic division and a few cholinergic junctions in the sympathetic division. Muscarinic receptors are G protein-coupled receptors, and their stimulation and subsequent activation of G proteins produce longer-lasting effects than nicotinic receptors. The response elicited can be either excitatory or inhibitory, depending on the activation or inactivation of specific enzymes.

 Core: Anatomy and Physiology

Chemical Signaling in the Endocrine System

JoVE 14974

A signaling cascade is a series of events that facilitates the transmission of information within or between cells, culminating in a targeted response in the recipient cell. As chemical messengers, hormones are pivotal in initiating and modulating these intricate signaling cascades based on their solubility.

Lipid-soluble hormones, such as steroid hormones, demonstrate an intracellular action. These hormones traverse cell membranes due to their lipid nature. Once inside the target cell, they bind to intracellular receptors in the cytoplasm or nucleus. For instance, testosterone, a steroid hormone, diffuses across the cell membrane into the extracellular fluid and binds to transport proteins in the bloodstream. Upon reaching the target cell, it dissociates and attaches to an androgen receptor inside the cell. This hormone-receptor complex enters the nucleus, regulating gene expression by binding to specific DNA sequences called hormone response elements. This process influences mRNA transcription and subsequent protein translation, altering cell metabolism and proliferation.

Conversely, water-soluble hormones, like peptides and amines, follow an extracellular mechanism. These hormones bind to receptors on the cell surface, initiating a cascade of events. The binding triggers intracellular signaling pathways, often involving second messengers like cyclic AMP. This cascade leads to cellular responses such as enzyme activation, altered membrane permeability, or gene transcription. Notably, the signaling pathway amplifies, allowing a small hormone concentration to elicit a robust cellular response. In summary, the divergent mechanisms of hormone action highlight the intricacies of cellular signaling pathways, ultimately shaping the physiological responses crucial for maintaining homeostasis.

 Core: Anatomy and Physiology

Gonadal and Placental Hormones

JoVE 14990

The gonads, namely the testes in males and the ovaries in females, are pivotal in producing gonadal hormones that orchestrate the intricate processes of sexual development and reproduction.

In males, testosterone is the primary gonadal androgen. It plays a central role in the maturation of male reproductive organs — the penis and testes. Additionally, testosterone is instrumental in the development of secondary sexual characteristics — a deep voice as well as facial and pubic hair growth — and the production of sperm.

On the female front, estrogen and progesterone are vital gonadal hormones. Estrogen is paramount for the development of female reproductive organs like the uterus and vagina. It also regulates secondary sexual characteristics, including breast development, and plays a crucial role in the menstrual cycle and the maintenance of pregnancy. Progesterone, on the other hand, prepares the uterus for the implantation of the zygote and ensures pregnancy by inhibiting uterine contractions.

The placenta takes center stage during pregnancy, secreting hormones vital for fetal growth and development. Human chorionic gonadotropin (hCG), a hormone that placental cells produce, supports pregnancy, particularly during the first trimester. Another placental hormone, human placental lactogen (hPL), regulates the mother's metabolism and stimulates the secretion of breast milk.

This interplay of gonadal and placental hormones is essential for the proper progression of reproductive processes and the sustained well-being of both mother and fetus during pregnancy.

 Core: Anatomy and Physiology

Superposition Theorem

JoVE 15052

The superposition principle is a fundamental concept stating that in a linear circuit, the voltage across (or current through) an element can be determined by summing the individual contributions of each independent source acting in isolation. When dealing with linear circuits containing multiple independent sources, this principle serves as a valuable tool for analysis. To apply the superposition principle effectively, one should focus on a single independent source at a time while deactivating all others. This approach yields the output (voltage or current) resulting from the active source.

The cumulative effect of all active sources can then be determined by algebraically adding their individual contributions. This simplifies the circuit analysis process. Notably, dependent sources remain unaffected as they are governed by circuit variables. It is worth noting that utilizing the superposition principle can lead to increased analytical effort. For example, when dealing with a circuit featuring three independent sources, one must analyze three separate, simplified circuits, each representing the contribution of an individual source. Despite this potential drawback, the superposition principle remains a valuable technique for simplifying complex circuits by replacing voltage sources with short circuits and current sources with open circuits.

 Core: Electrical Engineering

Design Example: Resistive Touchscreen

JoVE 15075

A device engineer plays a crucial role in designing user interfaces for mobile devices. One such interface is the resistive touchscreen, which fundamentally consists of two metallic layers: a flexible upper layer and a rigid lower layer, separated by a narrow gap. The high resistance between these two layers is a key characteristic of this design.

When a user touches the screen, the two layers make contact at a specific point known as the touchpoint. This contact reduces the resistance between the layers, effectively changing the electrical properties of the touchscreen at that point.

However, to accurately determine the precise location of the touchpoint, the touchscreen needs to be simplified as a one-dimensional system. The top layer, characterized by its length, resistivity, and cross-sectional area, is conceptually divided into two parts at the touchpoint. The resistances of these two sections are proportional to their lengths.

Connecting a voltage source between the two ends of the top layer transforms the circuit into a voltage divider configuration. The voltage drop at the touchpoint depends on the resistances of the two sections of the top layer.

By substituting the resistance values into the voltage divider equation and solving it, a relationship can be derived between the voltage drop at the touchpoint and its position. This relationship allows the device to calculate the exact location of the touchpoint based on the change in voltage.

This means that each touchpoint corresponds to a distinct voltage, enabling the system to accurately pinpoint the location of the user's touch. This precision is critical for the functionality of the touchscreen interface, ensuring it responds accurately to the user's input.

 Core: Electrical Engineering

First-Order Circuits

JoVE 15092

First-order electrical circuits, which comprise resistors and a single energy storage element - either a capacitor or an inductor, are fundamental to many electronic systems. These circuits are governed by a first-order differential equation that describes the relationship between input and output signals.

One common example of a first-order circuit is the RC (resistor-capacitor) circuit. These circuits are used in relaxation oscillators such as neon lamp oscillator circuits. When voltage is applied to an RC circuit, the capacitor begins charging, and the lamp acts as an open circuit. As the capacitor charges up to the required voltage to ionize the neon gas inside the lamp, the lamp suddenly becomes a short circuit. This causes the capacitor to discharge, creating a flash of light. Once the capacitor discharges fully, the process repeats, leading to a continuous flashing effect.

The time interval between these flashes depends on the time constant of the circuit, which can be adjusted by tuning the resistance (R) and capacitance (C) values. By carefully choosing these values, the frequency of the flashes can be controlled.

In tube lights, a different type of first-order circuit, known as an RL (resistor-inductor) circuit, is utilized. The choke coil in the tube light serves as the inductor, while the inherent resistance of the wire functions as the resistor. Upon voltage application, the choke resists a sudden increase in current, generating an electromotive force (emf) that rises with the applied voltage. This emf ionizes the gas within the tube light, causing it to illuminate.

In an RL circuit, the time constant is defined as the inductance (L) over the resistance (R). This time constant plays a vital role in determining how quickly the circuit responds to changes in the input signal. The larger the time constant, the slower the circuit responds, and vice versa.

 Core: Electrical Engineering

Impedances and Admittance

JoVE 15109

In the realm of AC circuits, passive circuit elements like resistors, inductors, and capacitors take on a different character when characterized by phasor voltage and current. Their behavior is expressed through impedance, a vital concept in AC circuit analysis.

Impedance is a measure of resistance to sinusoidal current flow in an AC circuit. Unlike their behavior in DC circuits, where inductors appear as short circuits and capacitors as open circuits, the behavior of these components in AC circuits is frequency-dependent. At high frequencies, inductors act as open circuits, while capacitors become short circuits.

Impedance is a complex quantity with a real part denoting resistance and an imaginary part representing reactance. Reactance can be either positive or negative, indicating inductive impedance when current lags behind voltage and capacitive impedance when current leads voltage. Impedance can also be represented in polar form, highlighting its magnitude and phase angle.

Equation1

Equation2

The reciprocal of impedance is admittance, which is measured in Siemens (S). Admittance represents the ease with which current flows through a circuit. It comprises conductance (real part) and susceptance (imaginary part). Admittance, like impedance, is a valuable tool in AC circuit analysis, enabling engineers to understand and manipulate electrical circuits operating under sinusoidal conditions.

 Core: Electrical Engineering

Introduction to z Scores

JoVE 17357

A z score (or standardized value) is measured in units of the standard deviation. It indicates how many standard deviations the value x is above (to the right of) or below (to the left of) the mean, μ. Values of x that are larger than the mean have positive z scores, and values of x that are smaller than the mean have negative z scores. If x equals the mean, then x has a zero z score. It is important to note that the mean of the z scores is zero, and the standard deviation is one.

z scores help find outliers or unusual values in any data distribution. According to the range rule of thumb, outliers or unusual values have z scores less than -2 or greater than +2.

 Core: Analytical Chemistry

An Introduction to the Chick: Gallus gallus domesticus

JoVE 5153

The chicken embryo (Gallus gallus domesticus) is an extremely valuable model organism for research in developmental biology, in part because most of their development takes place within an egg that is incubated outside of the mother. As a result, early developmental stages can be accessed, visualized and manipulated by simply creating a small hole in the eggshell. Since billions of chickens are raised worldwide for meat and egg production, scientists can easily and economically acquire large numbers of fertilized eggs throughout the year. Furthermore, chickens share significant genetic conservation with humans, so the genetic mechanisms that have been found to regulate chicken development are also relevant to our own biology.

This video focuses on introducing the domesticated chicken as a scientific model. The discussion begins with a review of chicken phylogeny, revealing the features that make them amniotes, like other birds, reptiles, and mammals. Highlights from the millennia of chicken research will be presented, ranging from Aristotle’s postulates about the function of extra-embryonic membranes to more recent, Nobel-prize winning discoveries in neuroscience. Additionally, some current examples of studies performed in chicken embryos will be provided, such as in vivo tracking of cell movements during development and the recruitment of blood vessels to developing tumors (a process known as angiogenesis).

 Biology II

Histological Staining of Neural Tissue

JoVE 5206

In order to examine the cellular, structural and molecular layout of tissues and organs, researchers use a method known as histological staining. In this technique, a tissue of interest is preserved using chemical fixatives and sectioned, or cut into very thin slices. A variety of staining techniques are then applied to provide contrast to the visually uniform sections. In the study of neuroanatomy, histological techniques are frequently applied to visualize and study nervous system tissue.

This video focuses on histological staining techniques for neural tissue. An overview of common brain stains is provided, including those that specifically mark neuronal cell bodies, like Nissl stains, and those that selectively highlight myelinated axons, like the Luxol Fast blue stain. Immunohistological techniques, which take advantage of the specific interaction between antibodies and unique cellular proteins, are also discussed. Next, the preparation of brain samples for staining is described, including the basic steps for fixation, embedding, sectioning, and rehydration of the tissue. The presentation also provides a step-by-step procedure for immunohistological staining followed by a Nissl stain, in addition to practical applications of these techniques.

 Neuroscience

An Introduction to Molecular Developmental Biology

JoVE 5328

Molecular signals play a major role in the complex processes occurring during embryonic development. These signals regulate activities such as cell differentiation and migration, which contribute to the formation of specific cell types and structures. The use of molecular approaches allows researchers to investigate these physical and chemical mechanisms in detail.

This video will review a brief history of the study of molecular events during development. Next, key questions asked by molecular developmental biologists today will be reviewed, followed by a discussion of several prominent methods used to answer these questions, such as staining, explant culture, and live-cell imaging. Finally, we will look at some current applications of these techniques to the study of developmental biology.

 Developmental Biology

Fear Conditioning

JoVE 5417

Fear Conditioning is a type of learning in which an association is established between a negative unpleasant event and a harmless stimulus. This leads to a fear of the harmless stimulus. This process is largely mediated by the amygdala, which is a brain region involved in emotions and stress reactions. Fear conditioning can be utilized in several ways to understand different aspects of learning and memory.

This video presents an overview of the principles behind fear conditioning, discusses the equipment and a generalized procedure used for this type of experiment. Finally, we'll review some real world applications of fear conditioning in behavioral neuroscience research today.

 Behavioral Science

Rotary Evaporation to Remove Solvent

JoVE 5501

Source: Dr. Melanie Pribisko Yen and Grace Tang — California Institute of Technology

Rotary evaporation is a technique most commonly used in organic chemistry to remove a solvent from a higher-boiling point compound of interest. The rotary evaporator, or "rotovap", was invented in 1950 by the chemist Lyman C. Craig. The primary use of a rotovap is to dry and purify samples for downstream applications. Its speed and ability to handle large volumes of solvent make rotary evaporation a preferred method of solvent removal in many laboratories, especially in instances involving low boiling point solvents.

 Organic Chemistry

DNA Methylation Analysis

JoVE 5550

Methylation at CpG dinucleotides is a chemical modification of DNA hypothesized to play important roles in regulating gene expression. In particular, the methylation of clusters of methylation sites, called “CpG islands”, near promoters and other gene regulatory elements may contribute to the stable silencing of genes, for example, during epigenetic processes such as genomic imprinting and X-chromosome inactivation. At the same time, aberrant CpG methylation has been shown to be associated with cancer.

In this video, the biological functions and mechanisms of DNA methylation will be presented, along with various techniques used to identify methylation sites in the genome. We will then examine the steps of bisulfite analysis, one of the most commonly used methods for detecting DNA methylation, as well as several applications of this technique.

 Genetics

Invasion Assay Using 3D Matrices

JoVE 5645

The extracellular matrix (ECM) is a network of molecules that provide a structural framework for cells and tissues and helps facilitate intercellular communication. Three-dimensional cell culture techniques have been developed to more accurately model this extracellular environment for in vitro study. While many cell processes during migration through 3D matrices are similar to those required for movement across rigid 2D surfaces, including adherence, migration through ECM also requires cells to modulate and invade this polymeric-mesh of ECM.

In this video, we will present the structure and function of ECM and the basic mechanisms of how cells migrate through it. Then, we will examine the protocol of an assay for tube formation by endothelial cells, whose steps can be generalized to other experiments based on 3D matrices. We will finish by exploring several other biological questions that can be addressed using ECM invasion assays.

 Cell Biology

Metabolic Labeling

JoVE 5687

Metabolic labeling is used to probe the biochemical transformations and modifications that occur in a cell. This is accomplished by using chemical analogs that mimic the structure of natural biomolecules. Cells utilize analogs in their endogenous biochemical processes, producing compounds that are labeled. The label allows for the incorporation of detection and affinity tags, which can then be used to elucidate metabolic pathways using other biochemical analytical techniques, such as SDS-PAGE and NMR.

This video introduces the concepts of metabolic labeling and show two general procedures.  The first uses isotopic-labeling, to characterize the phosphorylation of a protein. The second covers a photoreactive labeling to characterize protein-protein interaction within a Also three applications of metabolic labeling are presented: labeling plant material, labeling RNA to measure kinetics and labeling glycans in developing embryos.

Metabolic labeling is used to investigate the machinery of a cell. This is accomplished using chemical analogs to probe the biochemical transformations and modifications that occur. This video will show the principles of metabolic labeling, typical isotopic and photoreactive labeling procedures, and some applications.

Metabolic labeling can be conducted using a number of strategies. Here we will describe isotopic, photoreactive, and bio-orthogonal labeling.

Isotopic labeling is performed using structural analogs that are chemically identical to their natural counterparts, but have uncommon isotopes incorporated into their structure. In this L-lysine analog the carbon and nitrogen atoms are replaced with carbon-13 and nitrogen-15. Cells cultured in the presence of isotopic analogs will incorporate them into their biochemical structures. Metabolites are collected from the cells and purified for analysis. Samples with stable isotopes are analyzed using techniques such as mass spectrometry or NMR spectroscopy. Samples with radioactive labels are analyzed using liquid scintillation counting and x-ray films, which will be demonstrated in the isotopic labeling protocol.

Photoreactive labels are functional groups incorporated into proteins, which are stable until exposed to ultraviolet light. The functional group forms a reactive radical, which binds to the nearest protein. A common example, L-photo-leucine, contains a diazirine ring, which is a photoreactive crosslinker. In contrast to isotopic labeling, there is some chemical dissimilarity between photo-reactive chemical analogs and their natural counterparts. Cells may preferentially incorporate natural compounds over analogs. Therefore, it is important to perform photo-reactive labeling in medium free of the compound being mimicked. Once exposed to ultraviolet radiation, the photo-reactive groups in a labeled protein become unstable and highly reactive, causing it to cross-link with interacting proteins, creating a protein complex. Cross-linked complexes, act as snapshots that can then be analyzed using SDS-PAGE and mass spectrometry methods. This provides insights into what reactions are occurring in the metabolic pathway by identifying reaction species and how they interact by determining binding sites.

Bioorthogonal labeling strategies utilize analogs with small functional groups that have little to no reactivity with natural biomolecules. For example, azides are small functional groups, whose reactivity is said to be orthogonal to biochemical reactions. In the Staudinger ligation, a phosphine group attacks the azido group. This yields a transition state that intramolecularly reacts with a nearby ester, resulting in an amine-bonded ligand. Bioorthogonal functional groups incorporated into biomolecules can be ligated with detection tags such as fluorescent functional groups, and affinity tags such as antigens.

Now that some concepts and strategies for metabolic labeling have been discussed, let's look at the process in the laboratory.

The first step in a metabolic labeling experiment is to collect the protein of interest. To do this, cells are grown on a plate, and an expression method is used to promote synthesis of the desired protein. In this example, leucine-rich repeat kinases, or LRRK, are expressed. Disodium phosphate, containing radioactive phosphorus-32, is used as the analog. Proper measures must be taken to protect against ionizing radiation. This includes setting up a work area, wearing proper protective equipment, and checking for radioactive contamination. Once safety measures have been taken, the medium containing the isotopic analogs is prepared. The medium from the culture is removed and, replaced with one containing the isotopic chemical analogs and then incubated. Following incubation, the cells are lysed. The lysate is collected and purified.

After purification, proteins are resolved using SDS-PAGE and then transferred to a PVDF membrane. Autoradiography is performed by exposing the membrane to x-ray film and measured using a phosphor imager. Western blotting is used to measure relative protein levels in the PVDF membrane. In this example the phosphorylation levels of leucine-rich repeat kinases synthesized in 293T cells were measured. The autoradiogram shows how much phosphorous was incorporated into the protein. Western blotting elucidates the levels of the LRRK proteins. Image analysis software is used to obtain quantitative data of phosphorylation levels of the proteins.

In this next procedure, photoreactive labeling is demonstrated. First, the cells are prepared and cultured. The photoreactive analog is added to the cells at the mid-log phase and incubated. In this procedure p-benzoylphenylalanine is used. Samples are collected over intervals and put on ice. The samples are then exposed to obtain snapshots of the biochemical pathways over time. The proteins of interest are then purified and resolved using SDS-PAGE.

A photoreactive labeling strategy was used to identify compounds that interact with the protein of interest. Immunodetection with Western blotting shows protein bands that indicate higher molecular weight proteins are present in the irradiated samples. These are from cross-linking due to protein-protein interaction occurring during the UV irradiation.

Now that we've reviewed metabolic labeling procedures, let's look at some of the ways the process is used.

Metabolic labeling concepts can be extended to multicellular organisms. Plants are grown in a sealed environment, rich in stable isotopes to produced labeled plant material. Carbon dioxide containing carbon-13 is added to the enclosure, while nitrogen-15 rich fertilizer is used. The resulting harvested plant material can help answer questions about carbon and nitrogen cycling from the ecosystem.

Labeling enables the separation of newly synthesized RNA from older RNA. By changing the initial concentration of analog, the kinetics of new RNA synthesis can be determined. The results show that the concentration of 4-thiouridine affects how much new RNA is transcribed. Additionally, incorporation rates of the label into RNA can be directly quantified with a spectrophotometer.

Using biorthogonal click chemistry, glycans in a zebra fish embryo can be labeled. The eggs are injected with a labeling compound that results in alkyne labels on the glycans. The glycans in the larvae are then ligated to a dye compound at the desired development stage. The glycans in the embryos are then imaged. Glycans produced at different time points can be identified by labeling using different colors at different stages of embryo development.

You've just watched JoVE's video on metabolic labeling. This video described the concepts behind metabolic labeling and their strategies, went over two general procedures, and covered some of the uses of the techniques.

Thanks for watching!

 Biochemistry

Histotypic Tissue Culture

JoVE 5787

Although two-dimensional tissue culture has been common for some time, cells behave more realistically in a three-dimensional culture, and more closely mimics native tissue. This video introduces histotypic tissue culture, where the growth and propagation of one cell line is done in an engineered three-dimensional matrix to reach high cell density. Here, we show the harvesting of cells from donor tissue, followed by cell culture on an engineered construct.

 Bioengineering

Solubility

JoVE 11689

Solution, Solubility, and Solubility Equilibrium

A solution is a homogeneous mixture composed of a solvent, the major component, and a solute, the minor component. The physical state of a solution—solid, liquid, or gas—is typically the same as that of the solvent. Solute concentrations are often described with qualitative terms such as dilute (of relatively low concentration) and concentrated (of relatively high concentration).

In a solution, the solute particles (molecules, atoms, and/or ions) are closely surrounded by solvent species and interact through attractive forces. This dissolution process is called solvation. When water is the solvent, the process is known as hydration. For the solvation, the solute–solvent interactions should be stronger than solute–solute and solvent–solvent interactions. Precipitation is the opposite of solvation and occurs due to strong solute–solute interactions.

Solubility is the measure of the maximum amount of solute that can be dissolved in a given quantity of solvent. Temperature, pressure, and molecular polarity are some of the important factors that affect solubility. Solubility equilibrium is established when the dissolution and precipitation of a solute species occur at equal rates.

Like Dissolves Like

To predict if a solute will be soluble in a given solvent, the rule of thumb is “like dissolves like.” Polar or ionic solutes dissolve in polar solvents due to resulting ion–dipole or dipole–dipole interactions with the solvent molecules. Such interaction will not be possible with a nonpolar solvent. Nonpolar solutes dissolve in nonpolar solvents through intermolecular dispersion forces.

Hydrophilic, Hydrophobic, and Amphiphilic Compounds

Water is a polar solvent. Solutes that are soluble in water are called ‘hydrophilic’ or ‘water-loving’. For example, when solid KCl is added to water, the positive (hydrogen) end of the polar water molecules is attracted to the negative chloride ions, and the negative (oxygen) ends of water are attracted to the positive potassium ions. The water molecules surround individual K+ and Cl ions, reducing the strong forces that bind the ions together and letting them move off into solution as solvated ions.

A solute that is insoluble in water is termed as ‘hydrophobic’ or ‘water-fearing’. Such solutes, like oil, are unable to form hydrogen bonds with the surrounding water molecules due to the stronger solute–solute interactions. As a result, the solute particles cluster together and remain undissolved.

Compounds that have both polar and nonpolar groups are called ‘amphipathic’ or ‘amphiphilic’. For example, soaps, which are the salts of fatty acids. They have a hydrophobic tail of nonpolar hydrocarbons and a polar hydrophilic head. The cleaning action of soaps and detergents can be explained in terms of the structures of the molecules involved. The hydrocarbon (nonpolar) end of a soap or detergent molecule dissolves in or is attracted to nonpolar substances, such as oil, grease, or dirt particles. The ionic end is attracted by water (polar). As a result, the soap or detergent molecules become oriented at the interface between the dirt particles and the water, so they act as a kind of bridge between two different types of matter, nonpolar and polar. As a consequence, dirt particles become suspended as colloidal particles and are readily washed away.

This text is adapted from OpenStax Chemistry 2e, Section 11.1: The Dissolution ProcessSection 11.3: Solubility, and Section 11.5: Colloids.

 Core: Organic Chemistry

Lewis Acids and Bases

JoVE 11734

This lesson delves into Lewis acids and bases in the context of the octet rule for electron-deficient compounds. Here, the concept is discussed, emphasizing the group 13 elements like boron or aluminium. Since group 13 elements possess three valence electrons, they form trivalent compounds with a sextet of electrons and a vacant orbital for the central atom. Consequently, these electron-deficient compounds accept electrons from other species to complete their octet in a chemical reaction. They are referred to as Lewis acids per the 'generalized theory of acids and bases' proposed by Gilbert N. Lewis.

Lewis's theory dealt with compounds that were not under the purview of Brønsted's definition. He proposed that the electron-deficient compounds act as a Lewis acid where their valence-shell octets are completed in a chemical reaction. Hence, a Lewis acid is the species that accepts a pair of electrons to form a new bond.

In contrast, a Lewis base is defined as the species that donates an electron pair. This is elucidated using the specific example of aluminum chloride reacting with ammonia to form a Lewis acid-base adduct. Here, the electron pair is transferred between the oppositely charged species to satisfy the octet. The Lewis acids and bases concept is further reiterated by the reaction between electron-deficient boron trifluoride and an electron-rich ammonia, as shown in Figure 1.

Figure1

Figure 1. The reaction between Boron trifluoride and ammonia

Here, a significant charge is developed on the species. As the boron center has an empty orbital that could accept an electron, it localizes a positive charge. In contrast, the nitrogen center in ammonia accumulates a negative charge due to the presence of a lone pair of electrons. Hence when they interact, the lone pair of electrons in the valence shell of nitrogen is transferred to the boron atom in BF3, indicated by the curved arrow above. Thus, the formal charges on boron and nitrogen are balanced, and as a result, the Lewis acid-base adduct possesses no net charge.

The Lewis theory provides an addendum to the Brønsted theory that only uses proton transfer to define acid-base reactions by incorporating the transfer of a lone pair. Therefore, while all Brønsted–Lowry acids are protic acids, Lewis acids can be protic or aprotic. This is delineated using the example of hydrochloric acid. HCl is an acid, according to the Brønsted-Lowry definition, given its ability to donate a proton. It is also a Lewis acid since its hydrogen atom loses the shared electrons to chlorine while simultaneously accepting the pair of electrons from ammonia.

 Core: Organic Chemistry

Preparation of Alcohols via Substitution Reactions

JoVE 11771

Overview

Alcohols can be synthesized from alkyl halides via nucleophilic substitution reactions. The highly polar carbon-halogen bond in the substrate makes halide a good leaving group.  The hydroxide ion or water can act as a nucleophile to take the place of halide and form an alcohol. The substitution reactions occur via two different reaction pathways, SN1 or SN2,  depending on the nature of carbon attached to the halide.

Primary alcohols are synthesized from primary alkyl halides, and the reaction proceeds via the SN2 mechanism. The nucleophile attacks the halogen-bearing carbon from the side opposite to the carbon-halogen bond. However, in the presence of a strong nucleophile, a competing elimination reaction occurs as well.

Figure1

Figure_1: Parallel reactions of 1-bromobutane into substitution products and elimination products (proton abstraction).

The synthesis of secondary alcohols from secondary alkyl halides via substitution reaction is not favored since a mixture of products is formed from the competing SN2 and E2 reaction routes.

Figure2

Figure_2: Parallel reactions of 2-bromo-3-methylbutane into substitution products and elimination products (proton abstraction).  

Tertiary alkyl halides undergo SN1 reaction with a weak base such as water to produce tertiary alcohols along with alkene as a minor product due to a competing E2 elimination reaction.

Figure3

Figure_3: Parallel reactions of tertiary alkyl halides to elimination and substitution products.

If a strong nucleophile like sodium hydroxide is used, the E1 reaction dominates over SN1.

The nature of the reactant determines the stereochemistry of the product formed. If the halogen in the alkyl halide is connected to a chiral carbon, the resulting alcohol is a mixture of two enantiomers.

Figure4

Figure_4: Substitution reaction over an asymmetric carbon to yield a racemic mixture of optically active alcohols as the product

 Core: Organic Chemistry

Formation of Higher-order Actin Filaments

JoVE 11801

The polymerization of G-actin monomers into filamentous F-actin is a multi-step process. Once the F-actins are formed, they can bundle together in different arrangements to form higher-order networks and regulate cellular functions. Common examples include the formation of lamellipodia and filopodia at the cell's leading edge by actin reorganization in a migrating cell. The microvilli on the brush border epithelial cells are also formed through the F-actin network.

The high-order actin networks with straight F-actins are formed either through their bundled arrangement or by cross-linking them into gel-like networks. The dendritic networks are formed with the help of the branched F-actins. Actin binding proteins like fimbrin, fascin, and alpha-actinin form different types of actin filament bundles. Contrastingly filamin protein help in cross-linking the actin filaments into a gel-like network.

Actin bundling proteins

Actin-bundling proteins can arrange F-actins in either parallel or anti-parallel linear arrays. The bundles can be loose or tight depending on cellular functional requirements and the bundle's accessory protein. Monomeric proteins like fimbrin have two actin-binding domains and tightly bind to parallelly arranged adjacent actin filaments. These bundles can be found in microvilli in the small intestine. Contrastingly α-actinin is a dimeric protein having one actin-binding domain on each monomer. A helical spacer separates these actin-binding domains to form loose bundles with anti-parallelly arranged F-actins.

Actin cross-linking protein

Filamin is an actin cross-linking protein having two long flexible forms with one actin-binding domain on each arm. This allows flexible movement of filamin to form perpendicularly arranged actin filaments into a gel-like network. Some small proteins like transgelin have been reported to form dense meshworks.

 Core: Cell Biology

Destabilization of Microtubules

JoVE 11910

The destabilization of microtubules can occur during different stages of the microtubule lifecycle, such as nucleation or elongation. It can take place at either end of the microtubule or in the microtubule lattices as a whole. The lifespan of individual microtubules within a cell varies according to the cell type and stage of the cell cycle. During interphase, the lifespan of the microtubule is about 30 minutes, while during cell division, it is about 15 minutes. In axonal microtubules of neurons or axonemes of cilia and flagella, the microtubules have a longer lifespan.

Factors that Influence Microtubule Destabilization

Temperature is one of the factors that affect microtubule destabilization. In vitro studies have demonstrated that microtubules disassemble faster at 4°C while rapidly reassembling at 37°C. The next factor is the critical concentration (Cc) of microtubules, which is the concentration of free αβ-tubulin heterodimers at which the net polymerization of the microtubule is zero. The Cc influences the assembly or disassembly of microtubules. The microtubule destabilization occurs at heterodimer concentration lower than that of the Cc. In the cell, the microtubules may undergo abrupt transitions from catastrophe to rescue and vice versa, depending on the Cc.

Some microtubule-associated proteins act as microtubule-destabilizing agents (MDAs). These proteins primarily reduce the concentration of free αβ-tubulin heterodimers, inhibiting their longitudinal interactions with the microtubule. These MDAs facilitate their detachment from the microtubule upon attaching to the tubulin dimers in the protofilaments. Stathmin and kinesin-13 are the two most widely studied destabilizer MAPs responsible for increasing the frequency and duration of microtubule catastrophe.

Microtubule Associated Proteins that Influence Microtubule Destabilization

Stathmin-1 or oncoprotein 18 (Op18) was first discovered as an oncoprotein, highly expressed in breast and ovarian cancers and leukemia. These proteins are regulated by phosphorylation, which inactivates them and prevents their binding to the tubulin subunits. Stathmin is known to play a role in cell death. Stathmin can bind to both free tubulin subunits and those present in the microtubule protofilaments. Binding with free tubulin dimers changes their conformation, preventing their further binding to the microtubule filaments. In protofilaments, these destabilizing agents cause bending of the filaments, facilitating the removal of tubulin dimers.

Kinesin-13, a non-motile member of the kinesin superfamily, has its conserved motor domain at the center instead of the N- or C-terminal like other kinesins. This destabilizer reduces affinity between the tubulin subunits within the microtubule. Katanin, another microtubule destabilizer, has two subunits that sever the longitudinal bonds between the protofilaments of a microtubule. During cell division, katanin detaches microtubules from the centrosomes causing rapid destabilization of the spindle fibers. Their role in destabilizing microtubules has also been observed during interphase in proliferating cells.

 Core: Cell Biology

Oxidation of Alcohols

JoVE 11930

In this lesson, the oxidation of alcohols is discussed in depth. The various reagents used for oxidation of primary and secondary alcohols are detailed, and their mechanism of action is provided.

The process of oxidation in a chemical reaction is observed in any of the three forms:

  • (i) loss of one or more electrons,
  • (ii) loss of hydrogen,
  • (iii) addition of oxygen.

Oxidation is the opposite process of reduction, and hence, as carbonyls are reduced to alcohols, alcohols are oxidized to carbonyls. However, the oxidation of alcohols to carbonyls is dictated by the number of hydrogens present on the α-carbon linked to the hydroxyl group in the starting alcohol. Accordingly, while primary alcohols can be oxidized to aldehydes and further carboxylic acids, the secondary alcohols can only be oxidized to their corresponding ketones. Since there are no α-protons, tertiary alcohols cannot be oxidized. During the oxidation, there is a corresponding increase in the oxidation state of the central species.

Reagents and Mechanism

A popular reagent is the Jones reagent, a chromium trioxide solution in aqueous sulfuric acid in the presence of acetone. The reaction proceeds via an intermediate chromate ester and a subsequent E2 pathway to produce the carbonyl species. However, while the Jones oxidation ends at a ketone for secondary alcohol, the oxidation is repeated for primary alcohol resulting in a carboxylic acid. The other popular reagent used for oxidation is potassium permanganate. Similar to the Jones reagent, potassium permanganate is also a strong oxidizing agent converting the primary alcohol to a carboxylic acid. Hence, when an aldehyde is desired, a selective reagent like pyridinium chlorochromate or PCC should be used.

One major drawback of using these reagents is that they involve the higher oxidation states of chromium that are toxic. Accordingly, greener alternatives like Swern oxidation and Dess-martin oxidation have been designed. They employ reagents such as oxalyl chloride, dimethyl sulfoxide (DMSO), triethylamine, and dichloromethane, which are relatively non-toxic to convert the primary alcohols into aldehydes and secondary alcohols into ketones. In their mechanism, the Swern oxidation advances via an alkyl-sulfonium compound, while the Dess–Martin oxidation proceeds via a periodinane intermediate.

 Core: Organic Chemistry

Role of Myosin in Cell Migration

JoVE 12257

Myosins are multimeric motor proteins involved in various cellular processes such as migration, adhesion, and proliferation. Myosin II is the most common type in animal cells, which binds and cross-links actin filaments.

Myosin II  is a hexamer comprising two heavy chains with globular heads and coiled-coil tails, two regulatory light chains, and two essential light chains. The ATPase sites on the myosin heads hydrolyze ATP, and the released phosphate generates the force for contraction. It is this myosin-driven contraction of actin filaments and actin bundles that directs cell migration.

Myosin Promotes Adhesion

During cell migration, the lamellipodium at the cell front develops new focal adhesions that produce traction forces on the substratum. Studies have demonstrated that myosin is essential for the maturation of these adhesions, though the exact mechanism is not well understood. One hypothesis is that the myosin contractile force changes the conformation of cytoskeletal linker proteins, such as talin. This change exposes cryptic binding sites to actin and other linker proteins, thereby strengthening the adhesion. A second hypothesis states that adhesion proteins, such as integrins bound to actin, are clustered by the myosin-driven bundling of actin filaments. Thus, myosin-driven clustering of adhesion molecules, and changes in linker proteins, aid in the maturation of focal adhesions during cell migration.

Myosin and Organelle Positioning

Directional cell migration requires the establishment of cell polarity to determine the cell front and rear. In response to external stimuli, intracellular signaling cascades target various proteins, including myosin, to establish this polarity. Myosin is essential for repositioning the organelles, including the nucleus, Golgi apparatus, and the centrosomes, in alignment with the cell's polarity. For example, myosin anchors actin filaments around the nucleus by interacting with the linker protein nesprin embedded in the nuclear membrane. Thus, myosin helps to reposition the nucleus within the actin network as the cell changes shape during migration.

 Core: Cell Biology

Radical Halogenation: Thermodynamics

JoVE 12896

The thermodynamic favorability of a reaction is determined by the change in Gibbs free energy (ΔG). ΔG has two components- enthalpy (ΔH) and entropy (ΔS). The entropy component is negligible for alkane halogenation because the number of reactants and product molecules are equal. In this case, the ΔG is governed only by the enthalpy component. The most crucial factor that determines ΔH is the strength of the bonds. ΔH can be determined by comparing the energy between bonds broken and bonds formed.

Based on the thermodynamics of the reaction, radical halogenation of alkanes has a different order of reactivity for fluorination, bromination, and iodination. The ΔH for radical iodination is positive (+55 kJ/mol), which suggests that the ΔG value is also positive for this reaction. Therefore,  iodination is thermodynamically unfavorable, and the reaction does not take place. On the other hand, the overall ΔH for the radical fluorination of methane is large and negative (-431 kJ/mol), making the reaction thermodynamically favorable but highly exothermic and not having any practical use. The ΔH value for chlorination and bromination is -104 kJ/mol and -33 kJ/mol, respectively, making these reactions thermodynamically favorable and practically feasible. The reaction rate comparison between chlorination and bromination shows that bromination is slower than chlorination. The rate-determining step for this reaction is the first propagation step or the hydrogen abstraction step. The first propagation step for chlorination reaction is exothermic, and the energy of activation is small, while for bromination, this step is endothermic, and the energy of activation is large, which explains why bromination is slower than chlorination.

 Core: Organic Chemistry

Feedback Regulation of Calcium Concentration

JoVE 13327

Calcium is an essential signaling molecule required for various cellular functions. Calcium pumps and ion channels on cell and organellar membranes, such as those on the endoplasmic reticulum (ER), regulate calcium concentrations inside the cell. They remain closed, keeping the cytosolic calcium levels low at a resting state.

Various transmembrane receptors, such as G protein-coupled receptors (GPCRs), elicit a response to extracellular signals by increasing cytosolic calcium. Activated GPCRs stimulate phospholipase C, which produces inositol-1,4,5 trisphosphate (IP3). IP3 binds and opens IP3 gated calcium channels of ER, increasing the cytosolic calcium. The cytosolic calcium levels quickly spike, spreading across the cell by opening nearby IP3 gated calcium channels and ryanodine calcium channels. The sudden rise of calcium ions triggers various cellular responses. One such example is oocyte fertilization. Once a  sperm enters the oocyte's zona pellucida and fuses with the plasma membrane, the sperm’s phospholipase stimulates the release of IP3 from the oocyte’s plasma membrane. IP3 opens ER calcium channels, releasing calcium ions and facilitating calcium-induced calcium release (CICR) from the adjacent channels. This forms a calcium wave that moves across the egg starting from the point of sperm entry. The sudden rise in calcium stimulates modification in the egg’s surface and restricts the entry of other sperm. Calcium influx also pushes the zygote towards its first mitotic division by activating cell cycle regulators such as the cyclin-dependent kinases.

As the calcium levels go higher in a cell,  the ER calcium channels close, inhibiting the further release of calcium ions.  This feedback interaction between calcium ions and IP3 gated calcium channels causes repeated rise and fall in the cytosolic calcium, thus generating calcium oscillations in the cell. The oscillating cytosolic calcium regulates repeated cellular actions, such as the secretion of luteinizing hormone (LH) by the cells of the pituitary gland at every ovulation.

LH controls ovulation and is critical for maintaining female fertility. LH releasing hormone (LHRH) binds GPCRs on the pituitary gland cells' membrane. It triggers the release of calcium ions from ER lumen and facilitates the exocytosis of LH-containing secretory vesicles, releasing LH to the cell exterior.

 Core: Cell Biology

Overview Of Cell Separation And Isolation

JoVE 13371

Cell separation was first achieved in 1964 by S. H. Seal, who separated large tumor cells from the smaller blood cells using filtration. Two years later, Pohl and Hawk performed experiments on how cells respond differently to a nonuniform electric field based on the cell type. Such observations were the inception of cell separation methods, which allow isolating a single cell type from a heterogeneous sample.

Commonly used cell separation techniques include filtration, density gradient centrifugation, affinity-based separation using antibodies, and flow cytometric cell sorting. Additional specialized techniques include aptamer technology, buoyancy-activated cell sorting, and laser capture microdissection. The choice of separation method is based on the characteristics of the cell type to be isolated, such as surface charge, cell size, density, morphology, physiology, and surface markers. Usually, one or more of these characteristics are employed to isolate the particular cell type efficiently.

Applications of Cell Separation and Isolation

Isolated cells can be grown in vitro to produce cell lines, which have numerous applications in areas like pharmacology, immunology, and stem cell therapy. For example, the in vitro effects of drugs can be analyzed on specific cell populations. Cell separation is also essential in screening for the appropriate B-cells for monoclonal antibody production. Single-cell analysis, such as studying gene expression patterns and epigenetic effects, also relies on the isolation of specific cells to study. Many oncological studies that contribute to our understanding of cancer cells require isolating particular tumor cells from the tissue. Thus, cell separation and isolation methods are used virtually in all major fields of modern biology.

Challenges of Cell Separation Methods

Cell separation techniques face a high noise-to-signal ratio, with a small number of specific cells (target) against a large number of varied components that form the tissue. Various factors affect the purity of the final isolate; hence, the protocols require standardization for each target cell type and application. For example, the efficiency of cells dissociating from the tissue decides the yield of isolated cells. Excess digestion increases the presence of dead cells in the isolate, whereas incomplete digestion results in cell type contamination. Cell type contamination can also occur in affinity-based cell separation. While labeling a particular cell type with specific antibodies, non-specific binding may occur, introducing different cells into  the preparation. Additionally, the population size of the target cell type in the tissue is also a limiting factor — low-abundance cells often need tedious protocols to yield sufficient numbers in the isolated cell populations.

 Core: Cell Biology

Proteomics

JoVE 13388

A proteome is the entire set of proteins that a cell type produces. We can study proteomes using the knowledge of genomes because genes code for mRNAs, and the mRNAs encode proteins. Although mRNA analysis is a step in the right direction, not all mRNAs are translated into proteins.

Proteomics is the study of proteomes' function. It involves the large-scale systematic study of the proteome to denote the protein complement expressed by a genome. Scientist Mark Wilkins coined the term proteomics describing it as the ‘PROTein complement expressed by a genOME.’

As proteomics often complements genomics, it is useful when scientists want to test their research hypotheses. Even though all the cells in a multicellular organism have the same set of genes, the set of proteins produced in different tissues is different because it is dependent on gene expression. Thus, the genome is constant, but the proteome varies and is dynamic within an organism. In addition, RNAs can be alternately spliced (cut and pasted to create novel combinations and novel proteins). Also, many proteins modify themselves after translation by processes such as proteolytic cleavage, phosphorylation, glycosylation, and ubiquitination. There are also protein-protein interactions, which complicate studying proteomes. Although the genome provides a blueprint, the final architecture depends on several factors that can change the progression of events that generate the proteome. Because of this, there are different types of proteomics, such as  expression, structural and functional, which help study various aspects of the proteins.

The ultimate goal of proteomics is to identify or compare the proteins expressed from a given genome under specific conditions, study the interactions between the proteins, and use the information to predict cell behavior or develop drug targets. Just as analyzing the genome requires basic DNA sequencing technique, proteomics requires techniques for protein analysis. The basic technique for protein analysis is mass spectrometry that identifies and determines a molecule's characteristics. Advances in spectrometry have allowed researchers to analyze very small protein samples. X-ray crystallography, for example, enables scientists to determine a protein crystal's three-dimensional structure at atomic resolution. Nuclear magnetic resonance uses atoms' magnetic properties to determine the protein's three-dimensional structure in an aqueous solution. Scientists have also used protein microarrays to study protein interactions. Large-scale adaptations of the basic two-hybrid screen have provided the basis for protein microarrays. Scientists use computer software to analyze the vast amount of data for proteomic analysis.

 Core: Cell Biology

Cryo-electron Microscopy

JoVE 13404

Conventional electron microscopy (EM) involves dehydration, fixation, and staining of biological samples, which distorts the native state of biological molecules and results in several artifacts. Also, the high-energy electron beam damages the sample and makes it difficult to obtain high-resolution images. These issues can be addressed using cryo-EM, which uses frozen samples and gentler electron beams. The technique was developed by Jacques Dubochet, Joachim Frank, and Richard Henderson, for which they won the Nobel Prize in Chemistry in 2017.

X-ray diffraction and nuclear magnetic resonance (NMR) techniques are also used for obtaining high-resolution structures of biomolecules or studying their conformational changes. However, to get the X-ray structure, a molecule needs to be crystallized, which may not be possible every time. Even if a molecule is crystallizable, the crystallization process can alter its biomolecular structure, which is no longer representative of the native structure. Cryo-EM does not require a crystallized sample, and also enables visualization of molecular movements or interactions of biomolecules in their native state. Similarly, NMR technique is only limited to relatively small soluble proteins and is not suited for the non-polar proteins embedded in cell membranes. On the other hand, cryo-EM enables the study of larger proteins, membrane-bound receptors, or biomolecular complexes.

Cryo-EM is extensively being used in biochemistry to study the structure, function, and interaction of biomolecules. The infectious agents, such as viruses, can be identified and studied using cryo-EM. For example, during the Zika virus outbreak in Brazil, cryo-EM was utilized to create a 3D image of the virus structure so that the potential anti-viral drug targets could be identified.

 Core: Cell Biology

Clinical Applications of Epidermal Stem Cells

JoVE 13468

Epidermal stem cells (EpiSCs) are mainly located at the basal layer of the epidermis. These cells repair minor injuries of the skin and replace dead skin cells. However, EpiSCs’ cannot heal severe wounds such as major burns or those from diabetes or hereditary disorders. In such cases, culturing the epidermal stem cells from the patient is possible and has yielded successful treatment options, such as laboratory-grown skin grafts. These grafts are synthesized using a patient’s own EpiSCs containing skin tissue which is taken out from an unaffected region. These grafts, called cultured epidermal autografts (CEA), have been used to treat severe burn wounds for more than three decades and have saved many lives.

The use of autografts or autologous skin grafts has some risks. For example, taking skin from an unaffected area itself creates a new wound. Proper healing of such secondary wounds along with the primary burn wound adds to the treatment challenge. For such patients, allografts, or healthy skin tissue from another person, are used as grafts for wound healing. Although allografts are immunologically rejected by the recipient patient within weeks, they can help wound healing by inducing endogenous regeneration of skin tissue.

As a wound heals, it may also form a scar—a mark that remains after wound closure. Scars are thicker than normal skin and lack sweat glands. Large scars, especially those around joints, may limit the movement of body parts. Hence, a scar is an undesirable event during wound healing. However, the risk of scar formation can be minimized using EpiSCs. This can be achieved through the use of EpiSC-derived exosomes. Exosomes are extracellular vesicles containing biomolecules such as proteins or RNAs. In rodents, the EpiSC-derived exosomes have been shown to transport microRNAs to the injured site, promote wound healing, and prevent scarring.

 Core: Cell Biology

EPS and iPS Cells in Disease Research

JoVE 13485

Embryonic and induced pluripotent stem cells are excellent models for disease research because of their ability to self-renew and differentiate into most cell types. Somatic cells from a patient are isolated and reprogrammed into induced pluripotent stem cells or iPSCs. These iPSCs are later differentiated into the desired cell type, which mirrors the diseased cell of the patient. In this way, disease models have been created for investigating diseases such as Down syndrome, type I diabetes, and spinal muscular atrophy.

Drugs are usually tested in animal models, but using animals for research is expensive and may raise ethical concerns. Additionally, animal models may not accurately mirror human physiology or the disease under investigation. Patient-derived iPSCs have the same genetic makeup as the patient’s disease-affected cells. When differentiated, these cells can replicate the physiology of the diseased tissue and show a nearly accurate disease phenotype.

Drugs for treating cardiovascular, neurodegenerative, and liver disorders are currently being tested using cells derived from iPSCs. The cells are used to screen drugs for their effects and possible toxicity at varying doses. This can be part of a pre-clinical study of a medicine before it is tested directly in humans in clinical trials.

 Core: Cell Biology

Role Of Notch Signalling In Intestinal Stem Cell Renewal

JoVE 14080

Notch signaling was first discovered in Drosophila melanogaster, where it is involved in cell lineage differentiation. Notch signaling regulates the maintenance and differentiation of intestinal stem cells or ISCs by controlling the expression of atonal homolog 1 or Atoh1. Atoh1 directs cells to differentiate into secretory cells.

Direct cell-to-cell contact is needed for the activation of Notch signaling. The signal is initiated when a notch ligand binds to a receptor on an adjacent cell, also known as trans-binding. In contrast, the signal is inhibited when the ligand binds to the same cell's receptor or cis-binding. During trans-activation, ADAM protease cleaves the extracellular region of the notch receptor, and γ-secretase cleaves the intracellular domain, also called the Notch intracellular domain (NICD), which translocates into the nucleus and activates the Hairy enhancer of split (Hes). Hes, in turn, inhibits Atoh1 and prevents Notch receptor cells from differentiating into secretory cells.

During tissue damage to the intestinal epithelium, the Paneth cells de-differentiate and replace the damaged cells via Notch signaling. Dysregulated Notch signals in the intestinal cells can cause various diseases. For example, mutations in Hes1 or Atoh1 can give rise to colorectal tumors. Reduced Notch signaling is linked to the rare inflammatory condition known as Celiac disease, which causes an allergic reaction when gluten is consumed. Therefore, targeting the Notch signaling pathway can be a therapeutic strategy for treating certain intestinal disorders.  

 Core: Cell Biology

Calibration Curves: Linear Least Squares

JoVE 14518

A calibration curve is a plot of the instrument's response against a series of known concentrations of a substance. This curve is used to set the instrument response levels, using the substance and its concentrations as standards. Alternatively, or additionally, an equation is fitted to the calibration curve plot and subsequently used to calculate the unknown concentrations of other samples reliably.

For data that follow a straight line, the standard method for fitting is the linear least-squares method. This method minimizes the sum of the squared differences between the predicted and actual values.

The linear least square method plots the data points with the concentration on the x-axis and the measured analytical response on the y-axis. The equation of the line that best fits these data points is 'y = mx + c.' Here, y is the instrument's signal, x is the analyte concentration, m is the slope of the line, and c is the y-intercept. Once the best-fit equation has been determined, unknown concentrations can be determined with this equation by solving for x.

 Core: Analytical Chemistry

Solubility Equilibria: Overview

JoVE 14534

When a substance such as sodium chloride is added to water, it dissolves, forming an aqueous solution. The extent of dissolution is called solubility. The process of dissolution can exist in equilibrium, just like other chemical processes. Solubility equilibria are also called precipitation equilibria because the process of solubility can be reversible. The reverse of the solubility process is called precipitation.

Solubility is important in biological and environmental processes. A notable example of biological solubility is the effect of foods on tooth enamel, which consists of the mineral hydroxyapatite. Eating foods rich in sugars produces organic acids that dissolve hydroxyapatite, leading to tooth decay. Another example is calcium oxalate, a sparingly soluble salt that, if not flushed out by drinking sufficient water, will precipitate and form kidney stones over time.

Solubility equilibria follow Le Chatelier's principle, which states that if any force is applied to a reaction at equilibrium, the net reaction shifts towards whichever direction helps mitigate the stress from this force. Factors affecting the solubility equilibria of a sparingly soluble salt include temperature, solvent, common ion effect, pH of the solution, and the extent of complex ion formation. 

When a sparingly or moderately soluble salt is added to a solvent or solution, a dynamic equilibrium exists between the dissociated ions and the solid compound in the saturated solution. The equilibrium constant of a sparingly soluble salt is the solubility product constant or solubility product, Ksp, which is independent of the concentration of the solid form of salt because the concentration of the solid in a saturated solution is constant.

 Core: Analytical Chemistry

EDTA: Conditional Formation Constant

JoVE 14573

Each EDTA molecule has six binding sites: four carboxyl groups and two amino groups. The fully protonated form of EDTA is represented as H6Y2+. However, it can exist in different forms, H5Y+, H4Y, H3Y, H2Y2, and HY3, depending on the pH of the solution. In very basic solutions with pH > 10.17, the fully deprotonated form, Y4, is the predominant species that readily complexes with metal ions in a 1:1 ratio.

For the equilibrium reaction of the metal with the Y4 form of EDTA, the formation constant, Kf, is also known as the stability constant. Note that the Kf can be defined for any of the forms of EDTA in the solution. The value of Kf increases as the positive charge on metal ions increases. At a fixed pH, the different forms of the EDTA are in equilibrium so that the total molar concentration of all forms of EDTA equals the total molar concentration of the uncomplexed EDTA. The term cºEDTA is the total concentration of all forms of uncomplexed EDTA. At a fixed pH, the fraction of total EDTA present as Y4becomes constant, gives an expression where Kf′ is called the conditional formation constant or effective formation constant, which has a fixed value at a given pH. Therefore, it can calculate the equilibrium concentrations of free metal ions and the metal–EDTA complexes.

 Core: Analytical Chemistry

Mass Spectrum: Interpretation

JoVE 14589

An unknown compound can be established by identifying the molecular ion peak in the mass spectrum. The molecular ion peak is often weak or absent due to the predominance of fragmentation in high-energy electron beams. In such cases, a low-energy electron beam can be used to scan the spectrum to enhance the intensity of the molecular ion peak. Additionally, chemical ionization, field ionization, and desorption ionization spectra are used to obtain a relatively intense molecular ion peak.

To determine the empirical formula from the molecular ion mass, a high-resolution mass spectrometer like a double-focusing mass spectrometer, a time-of-flight mass spectrometer, an orbitrap mass spectrometer, or a Fourier transform ion cyclotron resonance mass spectrometer is required.

Comparing the intensities of isotope peaks (M+1 and M+2), which arise from the natural abundances of higher-mass isotopes, with the molecular ion peak (M+) gives a fair idea of the empirical formula of an unknown compound.

The number of carbon atoms present in the unknown compound can be deduced by comparing the M+1 to M+ peak intensity ratio to the expected M+1 contribution per carbon atom.

The nitrogen rule, stating that a molecule having an even molecular weight contains zero or even number of nitrogen atoms, predicts the number of nitrogen atoms present in the unknown compound.

Further, the degree of unsaturation predicted from the formula

Figure1

indicates the presence of multiple bonds or rings in the unknown compound. Lastly, the study of fragmentation patterns for compounds based on a series of guidelines helps identify signature fragments at characteristic m/z values, suggesting the presence of certain structural features.

 Core: Analytical Chemistry

Extraction: Advanced Methods

JoVE 14622

Metal ions can be separated from one another by complexation with organic ligands–the chelating agent– to form uncharged chelates. Here, the chelating agent must contain hydrophobic groups and behave as a weak acid, losing a proton to bind with the metal. Since most organic ligands used in this process are insoluble or undergo oxidation in the aqueous phase, the chelating agent is initially added to the organic phase and extracted into the aqueous phase. The metal-ligand complex is formed in the aqueous phase and is extracted back into the organic phase. During this process, the distribution of the metal-chelate complex between the organic and aqueous phases is independent of the initial metal concentration. It depends only on the pH of the aqueous phase and ligand concentration in the organic phase. For instance, consider the separation of divalent copper and lead. Aqueous copper(II) and lead(II) ions can be separated from each other by extracting with the mixture of dithizone in carbon tetrachloride. The plot of extraction efficiency versus pH reveals that copper can be quantitatively extracted into the organic phase if the pH of the aqueous phase is less than 5.5. After copper extraction, the pH of the aqueous phase can be increased to approximately 9.5, and lead ions can be extracted.

 Core: Analytical Chemistry

Disorders of the Skeletal Muscle

JoVE 14851

The clinical conditions affecting the skeletal muscle tissue are broadly categorized as musculoskeletal and neuromuscular disorders.

Musculoskeletal disorders

Musculoskeletal disorders involve injuries and conditions affecting the skeletal muscles and associated connective tissues. These disorders can arise from acute biomechanical stresses or chronic overuse and can occur across different age groups. Common injuries include sprains, fractures, and muscular strains, often resulting from accidents or intense physical activity. For example, a sudden slip or an improperly executed exercise can lead to muscle strains characterized by damage to muscle fibers. Additionally, repetitive movements can lead to conditions such as tendonitis, where the overuse of joints, like knees and wrists, causes inflammation in the tendons. Osteoarthritis, another prevalent condition, results from the long-term wear and tear of joint tissues and can significantly affect mobility and quality of life.

Neuromuscular disorders

Neuromuscular disorders primarily affect the neurons that control voluntary muscles or the interactions between muscles and neurons. These disorders often have a genetic cause and tend to affect most or all muscles in the body. One example of a neuromuscular disorder is muscular dystrophy, specifically Duchenne muscular dystrophy (DMD). DMD is a genetic disorder typically characterized by the mutation or absence of the dystrophin protein. Dystrophin plays a crucial role in maintaining the structural integrity of muscle cells. Dystrophin deficiency compromises the stability and strength of muscle fibers, leading to progressive muscle degeneration and weakness seen in individuals with DMD.

Another example of neuromuscular disorder is myasthenia gravis (MG), a chronic autoimmune neuromuscular disease causing varying degrees of skeletal muscle weakness. It is caused by a defect in transmitting nerve impulses to muscles. In MG, the immune system erroneously targets the synapses between the nerves and muscles, known as neuromuscular junctions. Specifically, the immune system produces antibodies that bind and obstruct the acetylcholine receptors on the muscle membrane of these junctions, thereby interrupting the transmission of signals from the nerve cells to the muscles. When these receptors are blocked or destroyed, muscle cells receive fewer nerve signals, resulting in muscle weakness. Typical symptoms of myasthenia gravis include drooping eyelids, difficulty swallowing, and general muscle weakness.

 Core: Anatomy and Physiology

Muscles of the Thorax

JoVE 14872

The thorax muscles are central to the body's respiration and provide essential support and movement for the upper body. They are intricately designed to facilitate the complex breathing process while also contributing to the structural integrity and mobility of the chest and upper limbs.

The diaphragm is at the core of thoracic musculature, the primary muscle involved in breathing. This expansive, dome-shaped muscle marks the division between the thoracic and abdominal cavities. It originates beneath the rib cage and lumbar vertebrae and attaches to the central tendon. When it contracts, the diaphragm moves downward and flattens, increasing the volume of the thoracic cavity and drawing air into the lungs. During exhalation, it relaxes and returns to its dome shape, helping expel air from the lungs.

 The rib cage has eleven pairs of intercostal muscles arranged into three groups: the external, the internal, and the innermost groups. The external intercostal muscles originate on the lower edge of the upper rib and attach to the upper edge of the underlying rib. They are most active during inhalation, where they contract to elevate the ribs and expand the chest cavity. On the other hand, the obliquely positioned internal intercostals have the opposite origin and insertion sites. They primarily assist in forced exhalation, pulling the ribs downward and inward to decrease the thoracic cavity's volume. The innermost intercostals, often considered a part of the internal layer, mirror the action of the internal intercostals and are similarly active during exhalation, helping compress the ribcage.

Additionally, the thorax houses muscles like the pectoralis major, pectoralis minor, serratus anterior, and subclavius, which facilitate shoulder and upper limb movements and stabilize the shoulder girdle. These muscles work with the respiratory muscles to support a range of upper body movements and postural adjustments.

 Core: Anatomy and Physiology

Nervous Tissue: Myelin

JoVE 14888

The myelin sheath is a multilayered lipid and protein covering that insulates the axon of a neuron, enhancing the speed of nerve impulse conduction. Axons without this sheath are referred to as unmyelinated. Two types of neuroglia, Schwann cells in the peripheral nervous system (PNS) and oligodendrocytes in the central nervous system (CNS) are responsible for producing myelin sheaths.

Schwann cells begin to form myelin sheaths around axons during fetal development. They wrap around a small portion of a single axon's length, creating multiple layers of glial plasma membrane around the axon. The inner part of this structure, comprised of the Schwann cell membrane, is the myelin sheath. The outer layer, which includes the Schwann cell's cytoplasm and nucleus, is called the neurolemma.

This outer layer is only found around axons in the PNS and is crucial in aiding axon regeneration after injury by forming a tube that guides and stimulates growth.

On the other hand, an oligodendrocyte in the CNS can myelinate parts of several axons. Unlike PNS, neurolemma is absent in the CNS because the oligodendrocyte cell body and nucleus do not envelop the axon. This absence, along with an inhibitory influence exerted by the oligodendrocytes, is believed to limit the regrowth of axons in the CNS following injury.

Myelination continues from birth to maturity, significantly increasing the speed of nerve impulse conduction. Consequently, infants, whose myelination is still in progress, exhibit slower and less coordinated responses to stimuli than older children or adults.

 Core: Anatomy and Physiology

Cerebrum: Anatomical Overview I

JoVE 14905

The main and largest component of the human brain is the cerebrum. The cerebrum consists of two main parts: the cerebral cortex, an outer layer with wrinkles or folds known as gyri and shallow grooves called sulci, and a deeper region beneath it. The cerebrum divides into two distinct hemispheres and contains five different lobes: the frontal, parietal, temporal, occipital, and insula. The central sulcus separates the frontal and parietal lobes and two functionally important gyri — the precentral and postcentral gyrus.

The lateral sulcus separates the temporal lobe from the frontal and parietal lobes, whereas the parieto-occipital sulcus separates the parietal and occipital lobes. The smallest lobe, the insula, is present deep in the lateral sulcus. The end of the occipital lobe and the separation of the cerebrum from the cerebellum are marked by a deeper sulcus known as the transverse cerebral fissure.

The cranial cavity protects the cerebrum perfectly, with the frontal lobes fitting in the anterior cranial fossa, the temporal and parietal lobes in the middle cranial fossa, and the occipital lobes present superior to the posterior cranial fossa.

 Core: Anatomy and Physiology

Cranial Nerves: Overview and Anatomy

JoVE 14921

The cranial nerves are an important part of the complex network of nerves in the human body. These nerves emerge directly from the brain and are responsible for transmitting essential information between the brain and various parts of the head and neck. There are 12 pairs of cranial nerves, systematically numbered using Roman numerals from I to XII, beginning from the anterior and moving to the posterior of the brain. Each cranial nerve is uniquely identified by names that reflect its function or the structure it innervates.

The cranial nerves are generally classified into three categories — sensory, motor, or mixed nerves. Whereas the sensory nerves detect information from the environment and the motor nerves control muscle movement, the mixed nerves have both sensory and motor functions.

Special Sensory Nerves: Nerves I (Olfactory), II (Optic), and VIII (Vestibulocochlear) fall under this category. They are responsible for the sensory modalities of olfaction, vision, hearing, and equilibrium, respectively. These nerves allow the human body to perceive and respond to external stimuli.

Motor Nerves: Nerves III (Oculomotor), IV (Trochlear), VI (Abducens), XI (Accessory), and XII (Hypoglossal) are classified as motor nerves. They primarily innervate muscles involved in eye, neck, and tongue movements, facilitating a range of actions from looking at different objects to swallowing and speech.

Mixed Nerves: Nerves V (Trigeminal), VII (Facial), IX (Glossopharyngeal), and X (Vagus) are known as mixed nerves because they carry both sensory and motor fibers. These nerves serve diverse functions, including facial sensation, mastication, taste, and salivation. They also control muscles in the oral cavity, pharynx, and some visceral organs. Their mixed nature allows them to participate in sensory perception and motor control of the head and neck regions and autonomic functions.

Parasympathetic Innervation

In addition to these functions, four cranial nerves—III (Oculomotor), VII (Facial), IX (Glossopharyngeal), and X (Vagus)—also play integral roles in the parasympathetic division of the autonomic nervous system (ANS). They are involved in involuntary functions such as regulating heart rate, digestion, and glandular secretions. For example, the Vagus nerve carries parasympathetic fibers to the heart, lungs, stomach, and digestive tract, and the Facial nerve carries parasympathetic fibers to stimulate salivary and lacrimal glands.

 Core: Anatomy and Physiology

Reflex Activity

JoVE 14939

A reflex activity is an automatic, involuntary response to specific stimuli. It is a part of our survival mechanism, designed to protect us from potential harm. For example, when a bright light suddenly shines into our eyes, we instinctively close them or look away. This is a simple reflex activity orchestrated by the nervous system without conscious thought or effort.

A reflex exam is a diagnostic procedure performed by a healthcare professional to evaluate the functionality of a patient's reflexes. The exam can also assess the patient's sensory and motor function, as well as their coordination and balance.

It involves the stimulation of specific sensory receptors, typically through tapping or striking certain areas of the body and observing the resulting involuntary muscle contractions or reflex responses. The exam focuses on assessing the deep tendon reflexes, such as the knee-jerk reflex, ankle reflex, or biceps reflex.

The plantar reflex test is an essential tool used in neurological exams to assess the integrity of the central nervous system. It involves stroking the sole and observing the resulting movement of the toes. In a normal response, the toes curl downward.

However, in some cases, an abnormal response may occur. This is referred to as the Babinski reflex. Instead of the normal downward flexion of the toes, the big toe extends upward while the other toes spread out. This kind of reflex is typical in infants but pathological in adults.

It indicates a neurological dysfunction, such as damage to the descending tracts of the somatic nervous system or damage to the primary motor cortex. 

 Core: Anatomy and Physiology

Disorders of the Autonomic Nervous System

JoVE 14959

The autonomic nervous system (ANS) is an intricate network of nerves that controls functions such as the regulation of heart rate, digestion, and blood pressure regulation. When this system malfunctions, it can lead to various disorders that affect multiple bodily functions. One common feature of many autonomic disorders is the involvement of smooth blood vessels, which play a crucial role in regulating blood flow throughout the body.

Raynaud's disease, also known as Raynaud's phenomenon, is a disorder characterized by episodes of reduced blood flow to the fingers and toes in response to cold temperatures or emotional stress. This condition is believed to arise from an overreaction of the ANS, specifically the sympathetic nervous system, which controls the constriction and dilation of blood vessels. In individuals with Raynaud's disease, the small arteries that supply blood to the extremities constrict excessively, leading to a decrease in blood flow. As a result, affected areas may turn white or blue and become cold or numb.

Autonomic dysreflexia is a potentially life-threatening condition that primarily affects individuals with spinal cord injuries above the T6 level. Typically, autonomic dysreflexia is triggered by a noxious stimulus, such as a full bladder or bowel impaction. In this disorder, the ANS responds inappropriately, leading to a sudden increase in blood pressure. This response is due to an imbalance between sympathetic and parasympathetic nervous system activity. While the sympathetic system causes vasoconstriction, the parasympathetic system counters this effect by promoting vasodilation. In autonomic dysreflexia, the exaggerated sympathetic response overrides the parasympathetic activity, resulting in uncontrolled vasoconstriction of the smooth blood vessels. This, in turn, leads to a severe increase in blood pressure that can have serious consequences if left untreated.

Hypertension, commonly known as high blood pressure, is a prevalent disorder that affects millions of people worldwide. While it can have various causes, dysregulation of the autonomic nervous system is often implicated in its development. The ANS plays a critical role in maintaining blood pressure within a normal range. The sympathetic nervous system, for instance, stimulates the constriction of blood vessels, leading to an increase in blood pressure, while the parasympathetic system promotes vasodilation, resulting in lowered blood pressure. In individuals with hypertension, there is often an elevated sympathetic response, leading to chronic vasoconstriction of the smooth blood vessels. The sustained elevation in blood pressure can have detrimental effects on the cardiovascular system, increasing the risk of heart disease, stroke, and other complications.

 Core: Anatomy and Physiology

Types of Hormones

JoVE 14975

Hormones are classified into four main groups: steroids, eicosanoids, amino acid-based derivatives, and peptide hormones.

Steroids and eicosanoids fall under the category of lipid-soluble hormones. Steroids are derived from cholesterol and feature four interconnected carbon rings with variable side chains. Notable examples include estradiol from ovaries and testosterone from testes, exemplifying the critical roles of these lipid-soluble hormones in reproductive physiology. Eicosanoids, derived from arachidonic acid, encompass prostaglandins and leukotrienes, exerting localized effects in various physiological processes like injury and infection.

In contrast, water-soluble hormones include amines, peptides, and glycoproteins. Amine hormones, synthesized from a single amino acid, include epinephrine, a tyrosine derivative commonly known as adrenaline. This hormone is pivotal in instigating the fight-or-flight response during stressful situations. Another notable example is melatonin, derived from tryptophan, which regulates the sleep-wake cycle. Peptide hormones comprise amino acid chains, such as the insulin from the pancreas. After meals, the insulin released in the blood signals cells to absorb excess glucose from the bloodstream, contributing to glucose homeostasis. Additionally, certain peptide hormones, like the thyroid-stimulating hormone, possess glycoprotein attributes featuring carbohydrate chains appended to the peptide structure.

 Core: Anatomy and Physiology

Cells and Secretions of the Pancreas

JoVE 14991

The pancreas, a vital organ within the abdominal cavity, plays dual roles in the digestive and endocrine systems, collaborating with exocrine and endocrine cells to maintain optimal digestion and blood sugar levels.

Exocrine function is carried out by acinar cells, organized into clusters known as acini. These cells contribute to digestion by releasing substantial quantities of enzyme-rich, alkaline digestive juices.

Concurrently, the dispersed clusters of endocrine cells throughout the pancreas are called the islets of Langerhans. Four primary types of endocrine cells are recognized for their roles in blood glucose regulation.

Alpha cells, responsible for producing glucagon, stimulate the liver to release stored glucose, elevating blood sugar levels. On the other hand, beta cells synthesize insulin, a hormone that facilitates glucose uptake by the body's cells, thereby reducing blood sugar levels.

Delta cells produce somatostatin, also known as the growth hormone-inhibiting hormone, which plays a regulatory role by inhibiting insulin and glucagon release. Another endocrine cell type, the pancreatic polypeptide (PP) cells, produces the pancreatic polypeptide, influencing the rate of nutrient uptake in the intestines.

This intricate interplay between exocrine and endocrine components ensures the pancreas' multifaceted contribution to digestion and blood sugar regulation, highlighting its essential role in maintaining overall metabolic balance.

 Core: Anatomy and Physiology

Thevinin's Theorem

JoVE 15053

Thévenin's theorem plays a pivotal role in electrical circuit analysis, offering a solution to the challenges posed by variable loads within a circuit. In practical applications, it is common to encounter circuits where certain elements remain fixed while others fluctuate, often referred to as the "load." A typical household electrical outlet serves as a prime example of a variable load, as it can be connected to a variety of appliances, each with its own unique electrical characteristics. Every time a new appliance is plugged in or removed, the entire circuit requires a fresh analysis. Thévenin's theorem provides a technique that replaces the fixed portion of the circuit with an equivalent circuit, thereby simplifying analysis and design.

Thévenin's theorem states that a linear two-terminal circuit can be represented by an equivalent circuit consisting of a voltage source in series with a resistor. This equivalent circuit effectively encapsulates the behavior of the original circuit from an external perspective. Two critical parameters define this Thévenin equivalent circuit: the Thévenin equivalent voltage (VTh), which corresponds to the open-circuit voltage at the terminals, and the Thévenin equivalent resistance (RTh), which represents the input or equivalent resistance at the terminals when all independent sources are deactivated.

To illustrate the concept further, consider two circuits shown in Figure 1 and Figure 2. These circuits are deemed equivalent when they share the same voltage-current relationship at their terminals. When the terminals are left open-circuited, resulting in no current flow, the voltage across the open terminals will be the Thévenin voltage. When all independent sources are turned off. Then the resistance seen across the terminals is the Thévenin resistance.

Figure1

Figure 1

Equation1

Figure2

Figure 2

Equation2

To practically apply this idea and determine the Thévenin resistance, we must consider two scenarios:

CASE 1: In situations where the network contains no dependent sources, all independent sources are deactivated and the input resistance is ascertained, which is the resistance seen between terminals a and b.

CASE 2: When the network incorporates dependent sources, continue to deactivate all independent sources. However, refrain from deactivating dependent sources since they are controlled by circuit variables.

It is worth noting that the value of Thévenin resistance (RTh) can sometimes be negative, signifying that the circuit is supplying power rather than consuming it. This scenario is feasible in circuits with dependent sources.

 Core: Electrical Engineering

Operational Amplifiers

JoVE 15076

The operational amplifier, often referred to as an op-amp, is a multifaceted building block of a circuit. This electronic component functions like a voltage-controlled voltage source and can also be used to create a voltage- or current-controlled current source. The design of an operational amplifier enables it to execute mathematical operations when external components like resistors and capacitors are linked to its terminals. An op-amp has the capacity to sum signals, amplify a signal, integrate it, or differentiate it.

The op-amps are popular in practical circuit designs due to their adaptability, affordability, and ease of use. An op-amp consists of a complex network of resistors, transistors, capacitors, and diodes. Op-amps are commercially available in different forms of integrated circuit packages. One common type is the eight-pin dual in-line package (DIP), as depicted in Figure 1.

Figure1

Figure 1: A typical op amp pin configuration

Pin 8 is not used, and pins 1 and 5 are of minor concern. The five significant terminals are the inverting input (pin 2), the noninverting input (pin 3), the output (pin 6), the positive power supply V (pin 7), and the negative power supply V (pin 4). The circuit symbol for the op-amp is a triangle, as shown in Figure 2.

Figure2

Figure 2: A typical op-amp circuit symbol

The op-amp has two inputs and one output, with the inputs marked with a minus (-) and a plus (+) to denote inverting and noninverting inputs, respectively.

An input applied to the noninverting terminal will appear with the same polarity as the output, while input applied to the inverting terminal will appear inverted as the output. As an active component, the op-amp necessitates a voltage supply for power. The current drawn from power supplies is a combination of the output current and the current needed to power the amplifier's internal circuitry.

 Core: Electrical Engineering

RC Circuit without Source

JoVE 15093

When a DC source is abruptly disconnected from an RC (Resistor-Capacitor) circuit, the circuit becomes source-free. Assuming that the capacitor was fully charged before the source was removed, its initial voltage, denoted as V0, can be considered as the initial energy that stimulates the circuit.

Applying Kirchhoff's current law at the top node of the circuit and substituting the current values across the components, a first-order differential equation is obtained. By rearranging the terms in this equation, integrating, and then taking the exponential on both sides, the natural response of the circuit is determined. The integration constant in this equation equals the initial voltage.

The voltage versus time graph shows that the initial voltage decays exponentially with time. This means that the charge on the capacitor gradually decreases, which in turn reduces the voltage across it.

The time constant of the circuit, represented by the Greek letter tau (τ), signifies the time required for the capacitor to discharge to 36.8% of its initial voltage. This time constant plays a critical role in determining the rate at which the capacitor discharges and, as a result, the speed at which the circuit responds to changes.

By substituting the value of tau into the voltage response expression, the current flowing through the resistor, as well as the power dissipated in the resistor, can be calculated. The power dissipated in the resistor is the rate at which energy is lost in the form of heat.

Integrating the power dissipated over time provides the total energy absorbed by the resistor. As time approaches infinity, this energy approaches the initial energy stored in the capacitor. This implies that the initial energy of the capacitor gradually dissipates in the resistor, eventually depleting the capacitor's charge.

In conclusion, understanding the behavior of RC circuits when the DC source is removed provides valuable insights into the transient response of these circuits. This knowledge is essential for designing and analyzing circuits in applications such as signal processing, power electronics, and communication systems, where the rapid charging and discharging of capacitors is a fundamental process.

 Core: Electrical Engineering

Node Analysis for AC Circuits

JoVE 15110

Consider an angioplasty system featuring a catheter equipped with a turbine, a critical tool for removing plaque deposits from coronary arteries. This intricate medical device operates using a circuit model reminiscent of a dual-node RLC circuit powered by a current-controlled voltage source.

To unravel the complexities of this system, nodal analysis is employed, a powerful technique founded on Kirchhoff's current law (KCL), which remains valid for phasors. AC circuits can effectively be analyzed using nodal analysis.

The process begins with gathering information about the input source voltage, inductance, and capacitance values. These data points can calculate the driving voltage for the catheter's shaft. Leveraging angular frequency, inductance, and capacitance values, the impedance across the inductor and capacitor is determined, mapping out a frequency domain circuit.

Equation1

Equation2

KCL and Ohm's law are applied at both nodes, yielding equations that describe the system's behavior. When simplified and integrated, these equations reveal that the shaft voltage precisely equals the source voltage.

This comprehensive analysis provides essential insights into the electrical operation of the angioplasty system. The voltage data can then be converted into the time domain, allowing for assessing and optimizing the system's performance for effective plaque removal in medical procedures.

 Core: Electrical Engineering

What is ANOVA?

JoVE 17358

The Analysis of Variance or ANOVA is a statistical test developed by Ronald Fisher in 1918. It is performed on three or more samples to check for equality between their means.

Before performing ANOVA, one must ensure that the samples used for this analysis have three crucial characteristics or statistical assumptions. The first assumption states that the samples should be drawn from normally distributed samples, while the second requires that all the drawn samples be randomly and independently selected. The third and last assumption states that the samples should be drawn from populations with equal variances.

There are two commonly used types of ANOVA: one-way ANOVA and two-way ANOVA. One-way ANOVA is used for the samples categorized by one factor, whereas two-way ANOVA is used when two factors categorize the samples.

Further, ANOVA is a helpful method that has broad practical applications. It can help a consumer select a washing machine or refrigerator after comparing different models or help a sociologist discern whether a person's income depends on their upbringing. ANOVA is used in environmental sciences to determine the variation in average pollution levels among several water bodies. As a result, ANOVA is widely applicable in fields such as life science, business administration, social science, forensic science, etc.

 Core: Analytical Chemistry

Basic Chick Care and Maintenance

JoVE 5154

Chicks (Gallus gallus domesticus) are a valuable research tool, not only for studying important concepts in vertebrate development, neuroscience, and tumor biology, but also as an efficient system in which to propagate viruses. Although eggs can be purchased from external suppliers and working with chicks requires very little specialized equipment, an understanding of proper handling procedures is required for normal embryo development.

This video will provide an overview of egg handling principles, including an explanation of the incubation parameters that can profoundly impact development: temperature, humidity, and egg rotation. Most experiments that use chicken eggs require access to the embryo within the shell, which is achieved by cutting a small, resealable hole, or “window.” This process is described in step-by-step detail, along with several other techniques essential for working with chicks, such as candling and India ink injection. Finally, the video will review some practical applications of these basic techniques in advanced scientific research.

 Biology II

An Introduction to Developmental Neurobiology

JoVE 5207

Developmental neuroscience is a field that explores how the nervous system is formed, from early embryonic stages through adulthood. Although it is known that neural progenitor cells follow predictable stages of proliferation, differentiation, migration, and maturation, the mechanisms controlling the progression through each stage are incompletely understood. Studying development is not only important for understanding how complex structures are assembled, but also for characterizing and treating developmental disorders. Since injury repair processes are similar to those that occur in development, this field is also a promising source of insight into when and how nervous system tissues regenerate.

This video provides a brief overview of the field of developmental neuroscience, including some key experiments that have advanced our understanding of the mechanisms controlling the formation of early neural tissue and the further specialization of those cells into discrete subsets of neurons. The discussion focuses on prominent questions that developmental biologists are asking and then demonstrates some of the methods that they use to investigate these questions. Finally, applications of the techniques are presented to provide insight into what it means to be a developmental neuroscientist today. The range of experiments demonstrated includes genetic manipulation of intact embryonic brains, targeted differentiation of stem cells into nervous system cells, and staining techniques that allow for the quantification of specific developmental events, like the formation of new connections between neurons.

 Neuroscience

Explant Culture for Developmental Studies

JoVE 5329

Explant culture is a technique in which living cells or tissues are removed from an embryo for continued development outside of the organism. This ex vivo approach allows researchers to manipulate and observe developing tissues in ways that are not possible in vivo. Once established, explant culture is frequently used to understand the role of genes and signaling molecules in organogenesis.

This video will first introduce the basic principles of explant culture and demonstrate a protocol to isolate and grow explanted mammalian tissues. Common genetic and molecular methods of manipulating explant cultures will then be discussed. Finally, the viewers will learn about how explant techniques are currently being applied to study organ development.

 Developmental Biology

An Introduction to Cognition

JoVE 5419

Cognition encompasses mental processes such as memory, perception, decision-making reasoning and language. Cognitive scientists are using a combination of behavioral and neuropsychological techniques to investigate the underlying neural substrates of cognition. They are interested in understanding how information is perceived, processed and how does it affect the final execution of behaviors. With this knowledge, researchers hope to develop new treatments for individuals with cognitive impairments.

JoVE's introduction to cognition reviews several components of this phenomenon, such as perception, attention, language comprehension, etc. Key questions in the field of cognition will be discussed along with specific methods currently being used to answer these questions. Finally, specific studies that investigate different aspects of cognition using tools like functional Magnetic Resonance Imaging (fMRI) or Transcranial magnetic stimulation (TMS) will be explained.

 Behavioral Science

Cyclic Voltammetry (CV)

JoVE 5502

Source: Laboratory of Dr. Kayla Green — Texas Christian University

A Cyclic Voltammetry (CV) experiment involves the scan of a range of potential voltages while measuring current. In the CV experiment, the potential of an immersed, stationary electrode is scanned from a predetermined starting potential to a final value (called the switching potential) and then the reverse scan is obtained. This gives a 'cyclic' sweep of potentials and the current vs. potential curve derived from the data is called a cyclic voltammogram. The first sweep is called the 'forward scan' and the return wave is called the 'reverse scan'. The potential extremes are termed the 'scan window'. The magnitude of reduction and oxidation currents and the shape of the voltammograms are highly dependent on analyte concentration, scan rates, and experimental conditions. By varying these factors, cyclic voltammetry can yield information regarding the stability of transition metal oxidation state in the complexed form, reversibility of electron transfer reactions, and information regarding reactivity. This video will explain the basic setup for a cyclic voltammetry experiment including analyte preparation and setting up the electrochemical cell. A simple cyclic voltammetry experiment will be presented.

 Analytical Chemistry

Chromatin Immunoprecipitation

JoVE 5551

Histones are proteins that help organize DNA in eukaryotic nuclei by serving as “scaffolds” around which DNA can be wrapped, forming a complex called “chromatin”. These proteins can be modified through the addition of chemical groups, and these changes affect gene expression. Researchers use a technique called chromatin immunoprecipitation (ChIP) to better understand which DNA regions associate with specific histone modifications or other gene regulatory proteins. Antibodies are used to isolate the protein of interest, and the bound DNA is extracted for analysis.

Here, JoVE presents the principles behind ChIP, discussing specific histone modifications and their relationship to gene expression and DNA organization. We then review how to perform a ChIP protocol, and explore the ways scientists are currently using this technique.

 Genetics

An Introduction to Endocytosis and Exocytosis

JoVE 5646

Cells can take in substances from the extracellular environment by endocytosis and actively release molecules into it by exocytosis. Such processes involve lipid membrane-bound sacs called vesicles. Knowledge of the molecular architecture and mechanisms of both is key to understanding normal cell physiology, as well as the disease states that arise when they become defective.

This video will first briefly review a few pivotal discoveries in the history of endo- and exocytosis research. Next, some key questions will be examined, followed by a discussion of the prominent methods used to investigate these problems, including cell labeling, fusion assays, and fluorescence imaging. Finally, it will explore current research being conducted by scientists in the field today.

 Cell Biology

Photometric Protein Determination

JoVE 5688

Measuring the concentration is a fundamental step of many biochemical assays. Photometric protein determination takes advantage of the fact that the more a sample contains light-absorbing substances, the less the light will transmit through it. Since the relationship between concentration and absorption is linear, this phenomenon can be used to measure the concentration in samples where it is unknown.

This video describes the basics of photometric protein determination and introduces the Bradford Assay and the Lowry Method. The procedure in the video will cover a typical Bradford assay. Applications covered include direct measurement of very small volumes of nucleic acids to characterize concentration and purity, determination of coupling efficiency of a biomimetic material, and another variation of photometric protein determination using Remazol dye.

Determining the concentration of a protein in samples is a fundamental step in many biochemical assays. Photometric determination can be done with small sample sizes. The more a sample contains light-absorbing substances, the less the light will transmit through it. This provides a quantitative measurement of the absorbing substances. These concepts are so fundamental to science that the articles that introduced two of the techniques are in the three most cited papers of all time. This video will show the concepts behind some of the most common photometric protein determination techniques, how they are performed, and how the gathered data is analyzed.

Photometric protein determination is based on the relationship between concentration and light absorbency. This is known as the Beer-Lambert Law, which states that the concentration of a light-absorbing species is proportional to its absorbance.

This principle underlies all photometric protein determination methods.

For direct absorption analysis, the absorbance values of unaltered protein samples are measured. Because of their aromatic side chains, tryptophan and tyrosine residues give the highest absorbance readings at a wavelength of 280 nm.

However, these amino acids-which are two of the least frequently found in proteins-are present in different amounts in every protein, so each determination is unique. To overcome this limitation, more complex assays-that are not dependent on these amino acids-were developed.

One example is the Bradford Assay, where colored dye is added to the sample. The dye, known as Coomassie Blue, responds proportionally-the more protein present, the more binding events with the dye.

Then, protein concentration is determined by measuring the absorbance of the bound Coomassie Blue dye, which absorbs light at 594 nm. However, the Bradford assay is linear over a short range of concentrations, so dilutions are often required before analysis.

The Lowry Method combines the Biuret reagent, an alkaline solution of copper ions that react with peptide bonds, and the Folin-Ciocâlteu reagent, which oxidizes aromatic protein residues. The resulting color change of the sample is proportional to the protein concentration.

The absorbance of the reduced Folin reagent can be determined at 750 nm. Like direct absorption, each protein has a unique response, and must be calibrated for the protein of interest. Now that we've reviewed the basic principles behind some of the most common assays, let's look at how direct absorption and the Bradford assay are performed.

To begin a direct absorption analysis, the spectrophotometer is calibrated with a blank to determine zero absorbance. Standard solutions are prepared for use in creating the calibration curve. Then, an aliquot of the first standard is added to a cuvette, and placed into the spectrophotometer.

The absorbance value at 280 nm is then recorded. This process is repeated for each standard, using a clean cuvette for each run. Once complete, a calibration curve is created by plotting the absorbance versus concentration. The slope of this line is the molar attenuation coefficient, which relates absorbance to concentration.

Next, the unknown sample is added to a cuvette, and the absorbance value is recorded. As the data analysis for the different photometric determination methods is similar, we will cover that after we look at the Bradford assay.

Here, the Bradford protein assay is performed with a BSA standard on a 96-well plate. To begin, BSA stock solutions are prepared.

The unknown solutions are diluted with deionized water to ensure that the concentrations are within the assay's range. Depending on the kit, the Coomassie dye may also require dilution. Then, the calibration curve is set up by adding the BSA standards to the 96-well plate.

Deionized water is added to reach the needed concentration to generate a standard curve. The unknown sample should be added to the plate in triplicates to ensure an accurate measurement is taken. Coomassie dye is next added to each well, mixing with the pipette.

Deionized water is added to an empty well as a blank, to measure the absorbance. After waiting 5 min for the dye to bind, the absorbance is measured in a plate-reader at 590 nm.

Now that we've performed a few assays, let's look at how to analyze the data. Each photometric protein determination method is based on the Beer-Lambert Law.

The measured absorbance of the standards is used to create a calibration curve, which is then used to determine the concentration of unknown samples. This curve can be manually plotted, though newer spectrophotometric tools will create the calibration curve once all standards have been measured. These systems will also calculate protein concentration as unknown samples are analyzed.

Now that we've reviewed how to analyze photometric protein determination data, let's look at some of the ways these procedures are utilized.

The principles of photometric protein determination can also be used to directly measure nucleic acid concentration. The nanodrop spectrophotometer accepts samples of very small volume onto an optically active pedestal. The absorbance is then measured, and the system automatically determines the nucleic acid concentration. Because proteins and other sources can interfere with measurements, sample purity is determined by analyzing the 260 to 280 nm and 260 to 230 nm absorbance ratios. Pure nucleic acids typically yield ratios of approximately 1.8 and approximately 2.0 for DNA and RNA, respectively.

Photometric protein determination can also be used in the production of biomimetic materials, which are inspired from nature to elicit specific cellular responses. Recombinant adhesins are bound to polystyrene beads to simulate bacterial attachment to host cells. The Bradford assay is used to determine the coupling efficiency of the recombinant adhesion to the beads in the production of the biomimetic material.

Alternative photometric protein assays can be used in the detection and characterization of protein antimicrobials. Remazol brilliant blue R dye is covalently bonded to heat-killed bacteria. The protein antimicrobial is incubated in the dyed solution. Then, the sample is centrifuged, and the absorbance of the supernatant at 595 nm is measured using a microplate spectrophotometer. Increased absorbance, by the soluble dye released into the supernatant from the labeled bacteria, is a quantitative measurement of enzymatic activity.

You've just watched JoVE's video on photometric protein determination. This video described the underlying principles of photometric determination, went over general procedures for some common assays, and covered some new advances in techniques. Thanks for watching!

 Biochemistry

Overview of BioMEM Devices

JoVE 5788

Bio-microelectromechanical systems, also called BioMEMs, are microscale devices that enable the use of small sample and reagent volumes for diagnostic devices in vivo and in vitro. These devices perform various functions such as filtration, sensing, or synthesis on the microscale, enabling cost savings and improved sensitivity.

This video introduces BioMEMs, touches on their use in the bioengineering field, and presents some prominent methods used in fabrication. Additionally, this video discusses some key challenges associated with miniaturization of devices, as well as some applications of the technology.

 Bioengineering

Introduction to Functional Groups

JoVE 11690

Functional groups are group of atoms with specific chemical properties that occur within organic molecules and sometimes denoted as “R”. Functional groups are found along the carbon backbone of macromolecules can form chains or rings of carbon atoms. Functional groups can “functionalize” a compound by enabling it to adopt different physical and chemical properties.  

Types of common functional groups

The table below summarizes some of the major functional groups in organic chemistry. (The functional groups are highlighted in red)

Class of organic compounds Skeletal Structures of Compounds and functional groups (red) Name of Compounds
Alkene Figure1 1-hexene
Alkyne Figure2 2-butyne
Alkyl halide Figure3 Chloromethane
Alcohol Figure4 Methanol
Thiol Figure5 Ethanethiol
Ether Figure6 Diethyl ether (ethoxyethane)
Sulfide Figure7 Dimethyl sulfide  (methylthiomethane)
Epoxide Figure8 Ethylene oxide (oxirane)

This text is adapted from Openstax, Chemistry 2e, Section 20: Organic Chemistry.

 Core: Organic Chemistry

Acid and Bases: Ka, pKa, and Relative Strengths

JoVE 11735

This lesson delves into a critical aspect of the relative strengths of acids and bases. The strength of an acid is evaluated by the acid dissociation into its conjugate base and a hydronium ion in water. The complete dissociation of a strong acid is confirmed with a very high concentration of hydronium ions. As a result, an incomplete dissociation process affirms a weak acid. Therefore, the equilibrium is in the forward direction for strong acids and backward for weak acids in these reactions.

Accordingly, the acid strength is defined by the concentration of undissociated acid molecules and hydronium ions. While the weak acid can be estimated via the equilibrium constant (Keq), it is constant for a dilute solution, and the change in water concentration is negligible. This observation leads to a modified equilibrium constant known as the acidity constant or dissociation constant, Ka. To define the acidity constant that is a scale of acidity, consider the generic acid-base reaction:

Figure1

Figure 1: Dissociation of a generic acid in water

Here, HA denotes the generic acid, and Adenotes its conjugate base. The following expression represents the acidity constant for this reaction.

Figure2

Figure 2: Acidity constant for a generic acid dissociation

This relationship focuses on the concentration of hydronium ions in the numerator. Accordingly, an increase in these ions leads to an increasing acidity constant and a stronger acid.  In organic acids, typically, the magnitude of Ka is spread across several orders. Hence, the strength of different acids is expressed in terms of pKa values, calculated as the negative logarithm of Ka:

Figure3

Figure 3: Expression of pKa

Here, the minus sign indicates the inverse relationship between the pKa value and acidity. As elucidated with benzoic acid versus hydrobromic acid, a higher pKa value equals a lower Ka value, which indicates a weaker acid.

By extension of the above principle, the pKa values can also establish the strength of a base. Since the dissociation of a base forms a conjugate acid, a stronger conjugate acid corresponds to a weaker base. For instance, consider methanol versus ethylamine. The conjugate acid of methanol with a pKa value of −3.8 is more acidic than the conjugate acid of ethylamine with a pKa value of 10.6. Hence, methanol is a weaker base than ethylamine.

 Core: Organic Chemistry

Regioselectivity of Electrophilic Additions-Peroxide Effect

JoVE 11772

In the presence of organic peroxides, the addition of hydrogen bromide to an alkene yields the isomer that is not predicted by Markovnikov’s rule. For example, the addition of hydrogen bromide to 2-methylpropene in the presence of peroxides gives 1-bromo-2-methylpropane. This addition reaction proceeds via a free radical mechanism, which reverses the regioselectivity. The free radical reaction mechanism involves three stages: initiation, propagation, and termination.

Figure1

In the first initiation step, an oxygen–oxygen bond in the radical initiator undergoes homolytic cleavage.

Figure2

The di-tert-butyl peroxide is an excellent free-radical initiator as the homolysis of the O–O bond requires just 159 kJ mol–1 (38 kcal mol–1) of energy.

Figure3

The second initiation step involves the exothermic (ΔH = –70 kJ mol–1) abstraction of hydrogen from HBr by the tert-butoxy radical. The abstraction of bromine, however, is thermodynamically unfavorable (ΔH = 163 kJ mol–1).

In propagation steps, a bromine radical reacts with an alkene to generate an alkyl radical. 

Figure4

The regioselective addition of bromine at the less substituted carbon in the presence of peroxide can be understood from the transition states. The transition state shows that the formation of the more substituted radical involves an attack by a bromine radical at the less substituted (and less hindered) carbon atom, which is lower in energy than the transition state for the less substituted radical. Another reason is the stability exhibited by the more substituted radicals owing to the hyperconjugation and inductive effect.

Figure5

Figure6

The reaction is terminated when radicals combine to yield non-radical products.

Figure7

Figure8

While the peroxide-mediated addition of HI to an alkene does not occur because the first propagation step is endothermic, the reaction with HCl does not proceed as the second propagation step is endothermic.

In the addition of hydrogen bromide to an alkene, the bromine radicals can attack the less substituted vinylic carbon from either face to an equal extent. Hence, when an alkene is stereogenic, a racemic mixture of products is obtained. 

Figure9

Figure10

 Core: Organic Chemistry

Overview of Myosin Structure and Function

JoVE 11802

Myosins are a family of molecular motor proteins, first identified in the skeletal muscles, where they are responsible for muscle contraction. Along with their role in muscle contraction, these proteins also play a role in the intracellular transport of molecules and vesicles. There are twenty-four classes of myosins based on their domain sequence and organization. Of the twenty-four, six classes (Myosin I, Myosin II, Myosin V, Myosin VI, Myosin VII, and Myosin X)  have been well characterized. The other roles of these molecular motor proteins include forming contractile rings during cytokinesis, organelle transport across polar actin filaments, aiding cell polarization, and signal transduction.

Myosin I is a monomeric protein with a globular head and a short tail. Unlike other myosin proteins, the tail domain can bind with a lipid membrane. The globular head with the actin-binding domain attaches to the F-actin. This myosin protein allows the intracellular transport of molecules and vesicles. Along with actin filaments, they are also present in intestinal cells as a component of the small cellular projections in microvilli.

Myosin II is the most widely studied member of this class of proteins, and its role in high-speed motility for muscle contraction has been well established. The general structure of myosin II is a dimer composed of two helically arranged polypeptide chains, each having a globular head, a narrow neck, and an alpha-helical tail. Each globular head has an actin-binding and ATP-binding domain where ATP hydrolysis occurs. The neck portion comprises a shorter light chain domain that acts as a lever for the neck movement, and the tail region has a long heavy chain domain that extends into the tail. The two tail polypeptide chains form a coiled-coil structure.

 Core: Cell Biology

Microtubule Associated Motor Proteins

JoVE 11911

Eukaryotic cells have different motor proteins for transporting various cargo within the cell. These motor proteins differ based on the filament they associate with, the direction they move within the cell, and the type of cargo they transport. Motor proteins that associate with microtubules are known as microtubule-associated motor proteins. There are two families of microtubule-associated motor proteins —Kinesins and Dyneins. Both these proteins assist in the transport of cellular cargos within the cell by hydrolyzing ATP molecules. The wide spectrum of cellular cargos transported includes organelles, vesicles, protein complexes, chromosomes, RNA-protein complexes, etc.

Kinesins

Kinesins are found in all eukaryotes.  The human genome has forty-five genes encoding kinesin proteins, Arabidopsis thaliana has sixty, while the budding yeast Saccharomyces cerevisiae has six. Kinesin-1 (conventional kinesin) was the first molecular motor protein discovered in a squid neuron. It is involved in transporting cellular cargo via the microtubules. These motor proteins have a highly conserved motor domain of ∼340 amino-acid residues at its N-terminal, with some exceptions. The kinesin superfamily is broadly divided into 14 families with distinct members based on their conserved motor domains. These proteins are plus end-directed, i.e., the transport of the organelles and vesicles occurs from the center of the cell towards the cell periphery, showing the centrifugal movement of cellular cargos. This motor protein superfamily members are also involved in microtubule destabilization, DNA repair, transcription regulation, mitotic spindle assembly, and cell signaling during cell growth regulation.

Dyneins

Dynein is a family of minus-end directed motor proteins, i.e., it transports cellular cargos from the cell periphery towards the center of the cell, showing centripetal movement. They were initially identified from the core of eukaryotic cilia and flagella. The core, also termed axoneme, is where the dyneins make the microtubules slide, resulting in the characteristic whip-like movement of cilia and flagella. These motor proteins are widely present in eukaryotes except for most flowering plants, which lack dynein motor proteins.

Dyneins are classified as axonemal and cytoplasmic based on their location and functions. Cytoplasmic dynein is responsible for intraflagellar transport, while axonemal dynein is involved in the locomotory function of cilia and flagella. During interphase, these motor proteins transport vesicles, organelles,  proteins, and mRNA particles, while in dividing cells, these proteins are responsible for spindle assembly.

 Core: Cell Biology

Preparation of Epoxides

JoVE 12094

Overview

Epoxides result from alkene oxidation, which can be achieved by a) air, b) peroxy acids, c) hypochlorous acids, and d) halohydrin cyclization.

Epoxidation with Peroxy Acids

Epoxidation of alkenes via oxidation with peroxy acids involves the conversion of a carbon–carbon double bond to an epoxide using the oxidizing agent meta-chloroperoxybenzoic acid, commonly known as MCPBA. Since the O–O bond of peroxy acids is very weak, the addition of electrophilic oxygen of peroxy acids to alkenes occurs with ease, thereby following syn addition. Hence, the epoxides are produced with the retention of the alkene configuration.

Epoxidation via Air Oxidation

Although peroxy-mediated epoxidation is the most common method for alkene oxidation, ethylene oxide is synthesized at the industrial scale via air oxidation by treating a mixture of ethylene and air in the presence of a silver catalyst.

Cyclization of Halohydrins

Cyclization of halohydrins of alkenes in the presence of a base also yields epoxides, and the reaction follows the SN2 substitution mechanism. Hence, the nucleophile—the oxygen anion—and the leaving group—the chloride ion—must orient anti to each other in the transition state to make the halohydrins cyclization feasible.

Figure1

In noncyclic halohydrins, this anti-relationship is achieved by an internal rotation. For instance, in 1-chloro-2-methyl-2-propanol, shown in Figure 1, the hydroxyl and the chloro group are not oriented anti to each other. To achieve the anti-relationship, the carbon-bearing chloro group undergoes an internal rotation, thereby making the nucleophile attack—from the backside of the C–X bond—and the ejection of the leaving group feasible. Thus, the epoxides formed via halohydrin cyclization also retain the alkene configuration.

Similarly, the cyclic halohydrins must undergo conformational changes to achieve the anti-relationship. For example, the halohydrin of cyclohexane, shown in Figure 2, undergoes a conformational change from diequatorial to diaxial to successfully form an epoxide.

Figure2

 Core: Organic Chemistry

Chemotaxis and Direction of Cell Migration

JoVE 12259

Cells can detect chemical cues in their environment and reorganize the cytoskeleton to migrate toward them or away from them. This directional migration, called chemotaxis, is essential during embryogenesis and development, immune response, tissue repair and regeneration, and reproduction. These chemical cues can either attract or repel the cell's movement. For example, axon development is determined by a combination of chemoattractants and chemorepellents that direct the growing axon towards the appropriate targets.

Sensing the Gradient

Cells exhibit chemotaxis across a chemical gradient by sensing the spatial difference in the chemical concentration between the two ends of the cell. Most eukaryotic cells can detect a difference of as little as 1 to 2 percent between the cell front and rear. Smaller cells, such as prokaryotes, cannot detect this spatial gradient because the distance between the front and rear of the cell is too small. Instead, they detect the gradient temporally by moving in random bouts and identifying the direction in which the chemical concentration increases. They then continue moving along that direction for a short distance before reverting to tumbling in random directions.

Chemotaxis in Cancer

Diseases involving abnormal cell migration usually exhibit abnormal functioning of cell surface receptors or unregulated expression of ligands for these receptors. During tumor metastasis, the presence of a specific receptor and its ligands determine the likely target organ for establishing secondary tumors. For example, the common secondary tumor sites in breast cancer patients are the lungs and bone marrow. The CXCR4 receptor found in breast cancer cells binds to the ligand SDF-1α, which is expressed in lung and bone marrow tissue. This ligand triggers directional migration of the metastatic cells towards these organs, thus establishing secondary tumors.

 Core: Cell Biology

Mass Spectrometers

JoVE 13038

This lesson details the instrumentation of a mass spectrometer—a physical instrument to perform mass spectrometry on analyte molecules and record the characteristic mass spectra. This is achieved via three chief functions:

  1. Conversion of the gas-phase analyte atoms/molecules into a beam of positive or negative charged ions by ionization.
  2. Separation of the charged species based on their mass-to-charge ratio.
  3. Recording the relative abundance of each type of ion.

In the ionization chamber of the mass spectrometer, the vaporized analyte in a vacuum is struck with high-energy electrons. The electron's energy is around 70 eV, sufficient to strip an electron from the analyte. The resultant molecular ion further fragments into charged species and neutral molecules. The mass spectrometer only records the mass of charged species, as the charge enables the control of molecules by an electric or magnetic field.

The molecular ion and its charged fragments are accelerated by a series of negatively charged accelerator plates positioned appropriately into the detector in the analyzing chamber. A magnetic field is applied on the path between the accelerator plate and detector, which causes a curve in the path of charged species. At a constant magnetic field, the radius of curvature depends on the molecular mass of the charged species. Placing a slit in front of the detector ensures that only charged species of a particular molecular weight reach the detector.

By varying the magnetic field, the charged species of all molecular weights can be recorded at the detector, each unique molecular weight species at a time. By scanning through all the magnetic fields, the mass spectrometer provides information on the relative abundance of all charged species as the mass spectrum.

 Core: Analytical Chemistry

Calmodulin-dependent Signaling

JoVE 13328

Calmodulin (CaM) is a calcium-binding protein in eukaryotes that controls various calcium-regulated cellular processes. It has four calcium-binding sites that bind calcium to form the calcium-calmodulin ( Ca2+-CaM) complex. GPCR stimulation increases the calcium levels in the cells that bind to CaM and induces a conformational change.

The Ca2+-CaM complex does not have enzymatic activity by itself. Instead, the complex binds downstream target proteins, including membrane proteins or enzymes, and activates them. For example, the Ca2+-CaM complex activates a family of protein kinases called Ca2+ calmodulin-dependent kinases (CaM kinases). CaM kinases phosphorylate target proteins such as transcription factors and alter their gene expression.

One such CaM kinases are CaM kinase II which is abundantly present in the nervous system. They consist of a stack of two giant rings, where each ring contains six copies of the enzyme. The enzyme has two domains, the kinase domain and the hub domain. A regulatory segment inhibits the kinase activity of the enzyme, keeping it in an inactive state.

The binding of the Ca2+-CaM complex to the regulatory segment unlocks the enzyme from an inactive to an active state. The Ca2+-CaM complex also activates nearby kinase domains that have popped out from its hub domain. The activated kinase domains of adjacent enzymes phosphorylate each other’s regulatory segments by autophosphorylation. Autophosphorylated CaM kinase II stays active even after the decay of the calcium signal. This allows the activated kinase to be a memory trace of a previous calcium spike and function as a memory device of the nervous system. Once protein phosphatase removes phosphates from the regulatory segment, CaM kinase II is switched off.

CaM kinase II can also decipher frequency changes in the calcium oscillations. At a low frequency of Ca2+ spikes, kinase becomes inactive as autophosphorylation cannot keep the enzyme active until the subsequent Ca2+ spikes. At a high frequency of Ca2+ spikes, the enzyme becomes only partially inactive between each Ca2+ spike. This leads to a progressive increase in the enzyme’s catalytic activity until it becomes maximally active when all of its domains are autophosphorylated.

 Core: Cell Biology

Cell Culture

JoVE 13372

Most vertebrate cells grow in vitro attached to a substrate as a monolayer, called adherent cultures. The flasks and plates used to grow cells are chemically treated to facilitate cell attachment. However, a few cell types, such as hematopoietic cells, can grow in a suspension. In contrast to adherent cultures, suspension cultures can grow in non-treated cultureware using magnetic stirrers or spinner flasks to agitate the culture media

Culture conditions

The growth medium is a crucial component for the optimal growth of cells. Previously, natural growth media like blood plasma, chicken embryo extracts, and amniotic fluid were used. In recent times, however, synthetic media such as Dulbecco's Modified Eagle Medium (DMEM) and Roswell Park Memorial Institute (RPMI) are preferred. These basal media contain amino acids, vitamins, inorganic salts, and a carbon source such as glucose. They are supplemented with fetal bovine serum (FBS) as a source of growth and adhesion factors, hormones, lipids, and minerals. Antibiotics are added to the growth media to prevent microbial contamination, and culturing protocols are performed under aseptic conditions.

The pH of the medium is usually maintained at 7.4 for mammalian cells. The growth medium also contains a buffering system such as sodium bicarbonate with exogenous CO2 for pH regulation. A pH indicator like phenol red is incorporated into the culture medium to indicate pH changes. An optimum temperature of 37℃ and a CO2 concentration of 5% are maintained in the incubator.

Plant cell culture

Like animal cells, plant cells can also be cultured in vitro. Plant tissue, known as explant, is cultured in a nutrient medium containing plant growth hormones, micronutrients, and a carbon source. Relative proportions of the plant growth hormones auxins and cytokinins in the medium decide the tissue type developed in the culture. A balanced ratio of these hormones results in a mass of undifferentiated cells called the callus. By altering the hormone ratio, the callus can be differentiated into root or shoot, thus generating a complete plant. Explants are usually cultured on a solid nutrient medium, while single-cell suspensions can be grown in liquid media. Growth conditions like temperature, light intensity, and photoperiod govern the plant growth in culture. Plant cell cultures are used to produce improved hybrid plants and conserve endangered plants.  Explant cultures are also used in the large-scale production of plant-derived products.

 Core: Cell Biology

Imaging Biological Samples with Optical Microscopy

JoVE 13389

Optical microscopy uses optic principles to provide detailed images of samples. Antonie van Leeuwenhoek designed the first compound optical microscope in the 17th century to visualize blood cells, bacteria, and yeast cells. In 1830, Joseph Jackson Lister created an essentially modern light microscope. The 20th century saw the development of microscopes with enhanced magnification and resolution.

In optical microscopy, the specimen to be viewed is placed on a glass slide and clipped on the stage (a platform). Once the slide is secured, the specimen on the slide is positioned over the light using the x-y mechanical stage knobs. These knobs move the slide on the surface of the stage but do not raise or lower the stage. Once the specimen is centered over the light, the stage position can be raised or lowered to focus the image. The coarse focusing knob is used for large-scale movements with 4⨯ and 10⨯ objective lenses; the fine focusing knob is used for small-scale movements, especially with 40⨯ or 100⨯ objective lenses.

When images are magnified, they become dimmer because there is less light per unit area of the image. Microscopes produce highly magnified images, therefore, require intense lighting. In a brightfield microscope, an illuminator provides this light, typically a high-intensity bulb below the stage. Light from the illuminator passes through the condenser lens (located below the stage), which focuses all of the light rays on the specimen to maximize illumination. The position of the condenser can be optimized using the attached condenser focus knob; once the optimal distance is established, the condenser should not be moved to adjust the brightness. If less-than-maximal light levels are needed, the amount of light striking the specimen can be easily adjusted by opening or closing a diaphragm between the condenser and the specimen. In some cases, brightness can also be adjusted using the rheostat, a dimmer switch that controls the intensity of the illuminator.

A brightfield microscope creates an image by directing light from the illuminator at the specimen; this light is differentially transmitted, absorbed, reflected, or refracted by different structures. The objective lens forms a real inverted image of the specimen. The final image visualized in the eyepiece or ocular is a magnified virtual image of the one formed by the objective.

Different colors can behave differently as they interact with chromophores (pigments that absorb and reflect particular wavelengths of light) in parts of the specimen. Often, chromophores are artificially added to the specimen using stains, increasing contrast and resolution. In general, structures in the specimen will appear darker, to various extents, than the bright background, creating maximally sharp images at magnifications up to about 1000⨯. Further magnification would create a larger image but without increased resolution. This allows us to see objects as small as bacteria, visible at about 400⨯ or so, but not smaller objects such as viruses.

At very high magnifications, resolution may be compromised when light passes through the small amount of air between the specimen and the lens. This is due to the significant difference between the refractive indices of air and glass; the air scatters the light rays before the lens can focus them. To solve this problem, a drop of oil can be used to fill the space between the specimen and an oil immersion lens, a special lens designed to be used with immersion oils. Since the oil has a refractive index very similar to that of glass, it increases the maximum angle at which light leaving the specimen can strike the lens. This increases the light collected and, thus, the image's resolution. A variety of oils can be used for different types of light.

This text is adapted from Openstax, Microbiology 2e, Section 2.3: Instruments of Microscopy.

 Core: Cell Biology

Electron Microscope Tomography and Single-particle Reconstruction

JoVE 13405

Transmission electron microscopy (TEM) can be used to determine the 3D structure of biological samples with the help of techniques such as electron microscope tomography and single-particle reconstruction. While single-particle reconstruction can examine macromolecules and macromolecular complexes in vitro conditions only, tomography permits the study of cell components or small cells in vivo.

Electron Tomography

Electron tomography can be performed either in TEM or STEM (scanning transmission electron microscopy) mode. STEM is a technique primarily used for obtaining tomograms of thick biological specimens. It combines the sample surface scanning methodology of scanning electron microscopy, such that the electrons transmitted or scattered at each point where the electron beam hits the sample are collected by a series of detectors.

A variation of the conventional electron tomography technique is dual-axis tomography, where the sample is tilted around two axes with respect to the electron beam. This results in two separate tomograms that are then aligned to generate a 3D image of the sample. Dual-axis tomography has two advantages over single-axis tomography: a better reconstruction of the sample's expanded features and increased sample depth resolution.

Another variation is the electron cryo-tomography, in which samples are imaged under cryogenic conditions because fixation and dehydration can damage the biological structures. The method is primarily used for thin samples, less than 500 nm in thickness, because thick samples block the electron beam. Therefore, the technique has been limited to purified macromolecular complexes, viruses, or small cells such as bacterial cells.

Single-particle Reconstruction

Single-particle analysis is typically used with cryo-electron microscopy to generate 3D structures with near-atomic resolution. It is most suited for large or dynamic macromolecular complexes that are difficult to crystallize, making it a substitute for X-ray crystallography. For complexes that can be crystallized, single-particle analysis is a method of choice to obtain high-resolution structures. The technique has been used for studying membrane proteins, protein complexes, chromatin structure, and macromolecular machines such as ribosomes, and proteasomes.

 Core: Cell Biology

Distinctive Features of Adult Stem Cells vs Cancer Stem Cells

JoVE 13469

A stem cell is an unspecialized cell that can divide without limit as needed and can, under specific conditions, differentiate into specialized cells.

Adult stem cells

Adult stem cells are tissue-specific; hence, they divide to develop the tissue from which they originate. One type of adult stem cell is the epithelial stem cell, which gives rise to the keratinocytes in the multiple layers of epithelial cells in the epidermis of the skin. Adult bone marrow has three distinct types of stem cells: hematopoietic stem cells, which give rise to red blood cells, white blood cells, and platelets; endothelial stem cells, which give rise to the endothelial cell types that line blood and lymph vessels; and mesenchymal stem cells, which give rise to the different types of muscle cells.

Cancer Stem cells

Cancer stem cells are a minor subpopulation of cells present in tumors and have the potential to renew and differentiate. For this reason, they are known to share similar characteristics with normal stem cells. In 1997, Bonnet and Dick were the first to find evidence of cancer stem cells by isolating a subpopulation of cells that expressed the surface marker CD34. It was found that the CD34 subpopulation had the potential to develop tumors in severe combined immunodeficiency (SCID) mice that were histologically similar to the donor.

While adult stem cells are generally tissue-specific, cancer stem cells multiply to produce cancer cells with the potential to terminally differentiate and a few cancer stem cells. Adult stem cells have limited self-renewal capacity, while cancer stem cells have indefinite self-renewal capability.  Adult stem cells are quiescent most of the time and have normal karyotypes. But the cancer stem cells are mitotically less active than cancer cells and often contain abnormal karyotypes.

 Core: Cell Biology

Phases of Wound Repair

JoVE 13557

Following injury, the integrity of the injured tissues must be reestablished. For example, in skin tissue, wound repair involves coordination among resident skin cells, blood mononuclear cells, extracellular matrix, growth factors, and cytokines to complete the healing cascade.

Formation of Blood Clot

In case of deep injuries, trauma to blood vessels results in blood loss. In the meantime, phospholipids released from the ruptured endothelial cellular membrane are converted into arachidonic acid and metabolites like thromboxane A2 and prostaglandin 2α. These factors promote vasoconstriction at the site of injury, generally lasting up to 5-10 minutes and resulting in brief hypoxia. Due to a lack of oxygen, cells and tissues surrounding the injury site shift ATP production via the anaerobic glycolysis pathway. The lactic acid produced at the end of anaerobic glycolysis reduces the pH in the adjoining tissues and cells. Blood vessel trauma and reduced pH significantly induce platelet activation, adhesion, and aggregation. Next, the blood clot is formed, sealing the injury site from external infection and establishing the temporary matrix composed of thrombin, collagen, fibronectin, and platelets. This matrix induces several cytokines and growth factors that are needed during the repair process.

Chemotaxis and Activation

Once the clot is formed, damaged cells at the injury site send a distress signal to immune cells in the body. This is followed by the recruitment of the neutrophils at the injury site. Prostaglandins E2 have a central role in the inflammatory response. They promote vasodilation and increase the blood flow to allow the movement of neutrophils. The neutrophils inhibit the growth of bacteria by releasing proteolytic enzymes. Even macrophages play a critical role in all the phases of wound repair, such as secretion of cytokines and growth factors like interleukin and tumor necrosis growth factor. They also promote fibroblast proliferation and angiogenesis at the wound site.

Extracellular Matrix Reorganization

Collagen is the major fibrous protein in the extracellular matrix (ECM) that imparts tensile strength and regulates cell adhesion to the tissues. The damage caused to ECM is restored in the remodeling phase of the wound repair. In the granulation tissue, ECM produced by fibroblasts is composed of type III collagen — a weaker structural protein. To cope with higher collagen demand, fibroblasts prefer to secrete type III collagen, and the rate of collagen production is highest. It is in the remodeling phase of the healing cascade that the matrix metalloproteinases (released by fibroblasts) remodel type III collagen into type I collagen that is stronger and has higher tensile strength. The arrangement of type I collagen into parallel bundles helps wound contraction and provides rigidity to newly formed tissues.

 Core: Cell Biology

Calibration Curves: Correlation Coefficient

JoVE 14519

In a linear calibration curve, there is a value called the calibration coefficient, denoted by 'r,' which measures the strength and the direction of association between two variables. The correlation coefficient value ranges from −1 to +1. A value of +1 indicates a perfect positive linear correlation, −1 denotes a perfect negative correlation, and 0 implies no correlation between the two variables. A positive correlation value establishes that as one variable increases, the other increases, and vice versa. On the contrary, a negative correlation value indicates that as one variable increases, the other variable decreases, and vice versa. Squaring the correlation coefficient results in the coefficient of determination, denoted by 'r2' or 'R2'. This value ranges from 0 to 1. A value closer to 1, such as 0.999, indicates an excellent fit, whereas a value close to 0 indicates a poor fit.

 Core: Analytical Chemistry

Solubility Equilibria: Ionic Product of Water

JoVE 14535

Pure water is a weak electrolyte; only a small amount ionizes into hydrogen and hydroxide ions. At any given temperature, the concentration of undissociated water is almost constant, so the ionic product of water is the product of the hydrogen and hydroxide ion concentrations, denoted as Kw. The square root of Kw gives the individual ion concentrations.

The ionic product of water varies with temperature, and its value is 1.0 x 10−14 at standard experimental conditions. Per Le Chatelier's principle, if the product of hydrogen and hydroxide ion concentrations becomes more than this value at any point, the excess ions will combine to form water molecules until it reaches equilibrium. Similarly, if the product of ionic concentrations falls below this value, water ionizes to form hydrogen and hydroxide ions to attain equilibrium. 

The nature of an aqueous solution–neutral, acidic or alkaline–is defined by the concentrations of the hydrogen and hydroxide ions. In the case of acidic or basic solutions, Kw is still the product of the concentrations of the hydronium and hydroxide ions, but these two concentrations will clearly not be equal to each other.

 Core: Analytical Chemistry

EDTA: Auxiliary Complexing Reagents

JoVE 14574

EDTA titrations are usually carried out in highly basic conditions, where the fully deprotonated form of EDTA, Y4, actively complexes with the free metal ions in the solution. Several metal ions precipitate as hydrous oxide (hydroxides, oxides, or oxyhydroxides) under these conditions, lowering the concentration of free metal ions in the solution. For this reason, auxiliary complexing agents or ligands such as ammonia, tartrate, citrate, or triethanolamine are used in EDTA titrations to prevent unwanted precipitation. These ligands bind strongly to the metal ions to form metal−ligand complexes that are less stable than the metal−EDTA complexes. The addition of EDTA during titration displaces ligands from the metal−ligand complexes, forming more stable metal−EDTA complexes.

 Core: Analytical Chemistry

Mass Analyzers: Overview

JoVE 14590

The mass analyzer is a crucial component of the mass spectrometer. In the ionization chamber, the vaporized sample is bombarded with a high-energy electron beam to generate a radical cation and further fragment into neutral molecules, radicals, and cations. A series of negatively charged accelerator plates accelerate the cations into the mass analyzer. The mass analyzer separates ions according to their mass-to-charge (m/z) ratios and then directs them to the detector. The common types of mass analyzers include magnetic sector, double-focusing, quadrupole, time-of-flight, and ion trap.

A magnetic field is applied in the magnetic sector analyzer, and the moving charges deflect to a curved path. For a given magnetic field, the radius of curvature depends on the m/z value, and the ions of a particular m/z value pass to the detector. As a result, the magnetic field of the instrument is continuously varied to sequentially pass ions in a broader m/z range to the detector.

A double-focusing mass analyzer uses a combination of the electrostatic sector and magnetic sector. As the ions exit the source, the electrostatic field focuses the ionic beam to minimize the kinetic energy spread and then leads the ions to the magnetic sector. The magnetic sector sorts the ions based on the deflection. This double-focusing of the beam helps in achieving high resolution.

The resolution of the mass spectrometer is defined as its ability to differentiate between two closely spaced mass spectra. A unit resolution indicates the ability of the instrument to differentiate the masses with a difference of one integer unit. For example, an instrument with unit resolution can distinguish m/z 50 from m/z 51, or m/z 100 from m/z 101, or m/z 500 from m/z 501. Generally, the resolution for commercial spectrometers is in the range of 500 to 500,000.

 Core: Analytical Chemistry

Dialysis

JoVE 14623

Dialysis is a diffusion-based purification process that separates analyte molecules from a complex matrix. This is accomplished by allowing molecules in the solution to pass through a semipermeable membrane into a liquid on the other side. The membrane is usually made of cellulose acetate or cellulose nitrate, and the second liquid must be miscible with the solution. Ions (e.g., chloride or sodium) or organic molecules (e.g., glucose) can pass through the membrane pores, which generally have diameters ranging from 1 to 5 nm. Larger molecules (e.g., proteins, hormones, enzymes) with diameters significantly greater than the pore diameter are retained in the original solution.

Dialysis is an inexpensive technique that is commonly used to purify proteins, hormones, and enzymes. During the procedure, a sample is placed in a dialysis bag, and both ends are sealed. The sealed bag is then immersed into a beaker filled with a solution called the dialysate, which has a composition different from the original solution. As dialysis proceeds, smaller molecules freely diffuse across the membrane, migrating from areas of high concentration to those of low concentration, until equilibrium is reached. Larger molecules that cannot enter the pores are left behind. Dialysis is complete once no net movement of the smaller molecules can occur.

Dialysis is a slow process, so it is often performed overnight or over several days. The driving force for diffusion is provided by the concentration gradient created by differences in analyte concentration across the dialysis membrane. The rate of molecular diffusion can be increased by raising the temperature and decreasing the thickness of the membrane. Because dialysis relies on equilibrium, it often needs to be repeated a few times to decrease the concentration of unwanted small molecules to an acceptable level.

 Core: Analytical Chemistry

Exercise and Muscle Performance

JoVE 14852

Exercise induces a range of adaptations in muscle tissue, depending on the type and duration of activity. Such physical training can be broadly categorized into two types: endurance exercises and resistance exercises.

Endurance exercises

Endurance exercises involve running, swimming, or cycling, which require repetitive movements with low force output. When a person engages in endurance exercise, a few noticeable changes occur in their skeletal muscles. For instance, the number of capillaries surrounding the muscle fibers increases. Additionally, the number of mitochondria within the muscle fibers also increases. Another change is that the fibers synthesize more myoglobin. While these changes occur in all fiber types, they are particularly significant in slow oxidative fibers that rely heavily on aerobic pathways. The result of these changes is that the muscles become more efficient in their metabolism, leading to greater endurance, strength, and resistance to fatigue. In fact, regular endurance exercise can even cause fast glycolytic fibers to transform into fast oxidative fibers. However, the transformation is not total or straightforward; the extent of these changes varies significantly between individuals and depends on factors like training intensity, duration, and genetic predisposition.

Resistance exercises

On the other hand, resistance exercises, such as weightlifting, sprinting, or jumping, require more strength than stamina. Such activities primarily engage fast muscle fibers that rely on ATP and glycogen for energy. Consistent resistance training enlarges muscle cells, eventually increasing strength. During high-intensity resistance exercise, muscles are pitted against high-resistance or immovable forces, and strength, rather than stamina, is the focus. Just a few minutes of high-intensity resistance exercise every other day can lead to a 50% increase in muscle mass within a year, even for someone not particularly strong. The additional muscle bulk primarily results from the increased size of individual muscle fibers, particularly the fast glycolytic variety, rather than an increase in the overall number of muscle fibers.

In summary, endurance and resistance exercises are essential for building physical fitness. While endurance exercises focus on improving stamina, resistance exercises aim to increase strength. Incorporating both activities into a workout routine can help individuals build muscle mass and improve their physical performance.

 Core: Anatomy and Physiology

Muscles of the Abdomen

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The abdominal wall encircles the abdominal cavity, providing flexible protection and shielding the internal organs from harm. It is bordered at the top by the xiphoid process and costal margins, at the back by the vertebral column, and at the bottom by the pelvic bones and inguinal ligament. The abdominal wall is divided into two regions — the anterolateral and posterior regions.

Anterolateral Region

The anterolateral region comprises five paired muscles classified into the lateral and anterior groups. The lateral group consists of three flat muscles — the external oblique, the internal oblique, and the transversus abdominis. The internal oblique muscles run perpendicular to the superficial external oblique muscles, while the transversus abdominis muscle mainly runs transversely around the abdominal wall. The aponeuroses of these muscles interlace and form the linea alba, a fibrous band of connective tissue running vertically from the xiphoid process of the sternum to the pubic symphysis.

The anterior muscles, specifically the rectus abdominis, and pyramidalis, are positioned bilaterally to the linea alba. The long, flat rectus abdominis muscle extends from the sternum and ribs to the pubic bone. Just beneath the lower portion of the rectus abdominis lies the pyramidalis, a small and often variable muscle that not everyone has. When present, it's located in front of the rectus abdominis and attaches to the pubic bone and linea alba.

Posterior Region

The posterior abdominal wall is located between the last thoracic and lumbar vertebrae and is made up of paired muscles, including the quadratus lumborum, psoas major, psoas minor, and iliacus. The quadratus lumborum is located deep in the lower back on either side of the lumbar spine. It is a key stabilizer muscle that extends from the iliac crest to the lower ribs and lumbar vertebrae, aiding in maintaining upright posture. The psoas major is a fusiform muscle that runs from the sides of the lumbar vertebrae to the lesser trochanter process of the femur. It works closely with the iliacus, a fan-shaped muscle that merges with the psoas major to form the iliopsoas muscle. Finally, the psoas minor, present in only about half of the population, runs parallel to the psoas major but is a much smaller and thinner muscle. When present, it assists in flexing the pelvis and spine.

 Core: Anatomy and Physiology

Electrochemical Gradient and Channel Proteins: An Overview

JoVE 14889

An electrochemical gradient is a fundamental concept in biology and chemistry. It regulates the movement of ions across cell membranes. This movement is influenced by two factors:

The electrical gradient: The electrical gradient across cell membranes refers to the difference in electric charge between the inside and outside of a cell.  This difference drives the movement of ions towards or away from the cells. For instance, if the inside of the cell is more negatively charged relative to the outside,  positively charged ions or cations are attracted to the inside of the cell, while negatively charged ions or anions are repelled.

The chemical gradient: Ions naturally move from areas of high concentration to areas of low concentration, a process called diffusion. This is driven by the principle that things tend to spread out and reach equilibrium.

Together, these gradients create an electrochemical gradient, which influences the movement of ions across cell membranes. It is crucial for various biological processes, like nerve signaling, muscle contraction, and nutrient uptake.

Channel Proteins

Channel proteins are specialized molecules embedded in cell membranes. They act like gates, allowing specific ions to pass through the membrane. These proteins are vital in maintaining the cell's internal environment and regulating processes such as nerve impulse conduction and muscle contractions.

Channel proteins have a specific three-dimensional structure that forms a pore or channel through the cell membrane's lipid bilayer. This channel allows ions to pass through but is selective, meaning it only allows ions of a specific size and charge to traverse. For example, potassium channel proteins are designed to allow only potassium ions (K+) to pass, while sodium channel proteins permit only sodium ions (Na+). This selectivity is determined by the shape and charge of the channel's interior.

Channel proteins can be regulated to control the flow of ions. Some channels open and close in response to changes in voltage (voltage-gated channels), while others respond to chemical signals (ligand-gated channels).

 Core: Anatomy and Physiology

Cerebrum: Anatomical Overview II

JoVE 14906

Each cerebral hemisphere can be divided into three main regions. The outermost region, the cerebral cortex, is a thin layer (2 to 4 millimeters thick) made up of gray matter, consisting of neuron cell bodies, dendrites, glial cells, and blood vessels. The middle region, or white matter, is primarily composed of myelinated nerve fibers organized into three types of large tracts: association fibers, commissures, and projection fibers. Association fibers connect different areas within the same cerebral hemisphere, while commissures, such as the corpus callosum, bridge the gap between the two hemispheres. Projection fibers link the cerebral cortex to vital structures like the thalamus, brainstem, and spinal cord.

Finally, deep within the white matter lies the innermost region, the basal nuclei, housing clusters of neuron cell bodies, which play a crucial role in motor functions. The first two nuclei, globus pallidus and putamen, are adjacent and together form the lentiform nucleus. The third nucleus is the caudate nucleus, which, when viewed from the side, resembles an arch with a large head and a slender tail-like structure above the diencephalon.

 Core: Anatomy and Physiology

Cranial Nerves: Types Part I

JoVE 14922

Cranial nerves are responsible for transmitting motor and sensory information between the brain and various parts of the body. There are twelve pairs of cranial nerves, with the first six being essential in sensory perception, motor control, and autonomic functions related to the head and neck.

Olfactory Nerve (Cranial Nerve I)

The olfactory nerve, or cranial nerve I, is unique as it is purely sensory and dedicated to the sense of smell. This nerve originates in the olfactory epithelium of the nasal cavity and extends to the olfactory bulb, which is part of the brain. The olfactory nerve's primary function is to transmit sensory data from the nasal cavity to the brain, allowing for the perception of odors. Unilateral or bilateral loss of smell (anosmia) can indicate damage to the olfactory nerve or bulbs.

Optic Nerve (Cranial Nerve II)

The optic nerve, or cranial nerve II, serves as the visual conduit between the eye and the brain. It is a sensory nerve that transmits visual signals from the retina to the visual cortex in the brain. The optic nerve is crucial for vision, and any damage to this nerve can result in visual impairments or blindness.

Oculomotor Nerve (Cranial Nerve III)

The oculomotor nerve, cranial nerve III, primarily has motor function. It controls most movements of the eye. It constricts the pupil and maintains an open eyelid through the levator palpebrae superioris muscle. The oculomotor nerve innervates several muscles, including the medial rectus, superior rectus, inferior rectus, and inferior oblique muscles of the eye. Damage to the nerve can result in ptosis (drooping eyelid), dilated pupil, and limited eye movement (e.g., inability to move the eye up, down, or medially).

Trochlear Nerve (Cranial Nerve IV)

The trochlear nerve is the smallest cranial nerve. It is also primarily a motor nerve and innervates the superior oblique muscle of the eye, which helps direct the gaze downward and rotate the eyeball. The trochlear nerve is unique because it is the only one that emerges dorsally from the brainstem and innervates a muscle on the contralateral side.

Trigeminal Nerve (Cranial Nerve V)

The fifth and largest cranial nerve is the trigeminal nerve. It has both sensory and motor functions. It is divided into three branches: the ophthalmic, mandibular, and maxillary nerves, which supply sensation to the face and motor functions to the muscles of mastication. The trigeminal nerve is essential for facial sensation, as well as biting and chewing movements. Damage to the nerve can result in facial numbness, pain (e.g., trigeminal neuralgia), or weakness in the muscles of mastication.

Abducens Nerve (Cranial Nerve VI)

The abducens nerve primarily has a motor function. It controls the lateral rectus muscle of the eye, which is responsible for abducting the eyeball. The abducens nerve is crucial in coordinating eye movements and ensuring that both eyes can move together in the same direction. A lesion in this nerve can cause difficulty moving the eye outward, leading to horizontal diplopia (double vision).

 Core: Anatomy and Physiology

Somatic Spinal Reflexes

JoVE 14940

Somatic spinal reflexes are rapid, involuntary muscular responses to external stimuli that involve the somatic musculature and the spinal cord.

One of the most well-known somatic spinal reflexes is the stretch reflex, which is activated by the sudden stretching of a muscle. This reflex involves the activation of specialized sensory receptors called muscle spindles, which are located in the muscle tissue and detect changes in the length and speed of muscle contractions. When a muscle is suddenly stretched, the muscle spindles activate and send a signal to the spinal cord, triggering a reflexive contraction of the same muscle. This is a monosynaptic reflex, meaning a single synapse between the sensory and motor neurons, without any interneurons in between.

Another important somatic spinal reflex is the tendon reflex, which is triggered by the tension of a muscle tendon. This reflex involves the activation of specialized sensory receptors called Golgi tendon organs, which are located in the tendons connecting muscles to bones. When a muscle tendon becomes tensed during muscle contraction, the Golgi tendon organs activate and send a signal to the spinal cord, triggering a reflexive relaxation of the same muscle. This is a polysynaptic reflex, meaning that the reflex arc comprises interneurons present between the sensory and motor neurons.

In addition to the stretch and tendon reflexes, the somatic spinal reflex system also includes the polysynaptic flexor reflex, which is responsible for withdrawing a limb from a painful or damaging stimulus, and the crossed extensor reflex, which is responsible for maintaining balance and stability during movements involving the arms and legs. The flexor reflex involves activating flexor muscles in response to a painful or damaging stimulus. In contrast, the crossed extensor reflex involves activating extensor muscles on the opposite side of the body to maintain stability during the movement.

 Core: Anatomy and Physiology

Introduction to Special Senses

JoVE 14960

Sensory receptors play an integral part in comprehending our external and internal environments. They receive diverse stimuli, converting them into the nervous system's electrochemical signals. This conversion occurs as the stimulus alters the sensory neuron's cell membrane potential, instigating the generation of an action potential. This action potential is subsequently transmitted to the central nervous system (CNS), which integrates with other sensory data or higher cognitive functions. This leads to a conscious awareness of the original stimulus. The central integration may eventually induce a motor response.

Terminology such as 'sensation' or 'perception' is employed intentionally when delineating sensory function. Sensation refers to the initiation of sensory receptor cells concomitant with stimulus exposure. Conversely, perception involves the brain's interpretation of sensory stimuli into recognizable patterns. While sensation is a prerequisite for perception, it does not guarantee the latter. Structures or cells that discern sensations are known as receptors. The transformation of a receptor cell occurs directly in reaction to a stimulus. A transmembrane protein receptor within the cell membrane instigates a physiological alteration in a neuron. This typically transpires via ion channel opening or modifications to cell signaling pathways. Activation of the transmembrane receptors occurs due to chemicals named ligands. For instance, specific molecules present in food may act as ligands for taste receptors. Though not strictly defined as receptors, other transmembrane proteins respond to mechanical or thermal variations. These physical transformations in the proteins can amplify ion translocation across the membrane and trigger an action or a graded potential in sensory neurons.

Sensory Receptors and Perception of Environmental Stimuli

Environmental cues instigate the activation of specialized sensory cells within the peripheral nervous system. These sensory cells are diverse, each specifically attuned to varying forms of stimuli. These sensory cells can be categorized on three primary bases: cellular morphology, spatial location relative to the sensed stimuli, and functional characteristics. Structurally, sensory cells can be categorized by their distinct cellular type and their spatial positioning close to the stimuli they are equipped to perceive. From a functional standpoint, these cells can be classified by their unique ability to transduce stimuli - the process by which a physical stimulus, light, or chemical alteration translates into a change in the cell membrane potential.

Structural Receptor Types

Sensory information translation is principally performed by

  1. Neurons possessing non-encapsulated nerve endings, where dendritic extensions are integrated within the tissue receiving sensory input
  2. Neurons with encapsulated nerve termini, where sensory nerve endings are enveloped in a connective tissue layer, augmenting their sensitivity
  3. Specialized receptor cells are constructed with explicit structural components to interpret a specific stimulus category.

An instance of neurons with non-encapsulated nerve endings is the nociceptive and thermoreceptive neurons within the skin's dermis. Pacinian corpuscles, neurons with encapsulated nerve termini sensitive to mechanical pressure and touch, are also found in the skin’s dermis. Photoreceptors in the retina exemplify a specialized receptor explicitly designed to respond to light stimuli.

Receptors can alternatively be stratified based on their proximity to the stimulus source. Exteroceptors represent a class of receptors near an external environmental stimulus, exemplified by the somatosensory receptors within the dermal layers. Interoceptors, in contrast, interpret stimuli originating from the internal body organs and tissues. An illustration of this would be receptors that register fluctuations in blood pressure within the aortic or carotid sinus. On the other hand, proprioceptors are located adjacent to moving body segments, such as muscles, and are critical for interpreting positional changes in the tissues during movement.

Receptor Functionality Classification

Receptors can be further bifurcated based on their mechanism of converting stimuli into membrane potential alterations. Generally, stimuli can be classified into three broad categories.

Certain stimuli comprise ions and macromolecules that interact with receptor proteins on the cell membrane, instigating a change when these biochemical entities diffuse across the cell membrane.

Other stimuli correspond to changes in physical environmental conditions under which receptor cell membrane potentials are susceptible to alterations.

The remaining category of stimuli encompasses electromagnetic radiation, notably visible light, which can be sensed by the human eye. Interestingly, several organisms possess specialized receptors absent in humans, such as snakes' heat-detecting sensors, bees' sensors for ultraviolet light, or birds' receptors for magnetic fields during migration.

Cells that interpret signals can be further classified into various types, depending on the nature of the stimulus they transduce. Chemoreceptors, for instance, process chemical signals, enabling the recognition of the taste or odor of an entity. Osmoreceptors detect fluctuations in body fluid solute concentrations. Nociceptors primarily serve as a chemical sensing system, discerning the existence of chemical substances arising from tissue harm or similar intense stimuli, thereby perceiving pain. Physical stimuli, encompassing pressure, vibration, auditory sensations, and positional awareness (equilibrium), are discerned via a mechanoreceptor. Thermoreceptors are specialized receptors that detect temperature changes, with distinct types responsive to temperatures either higher (heat) or lower (cold) than the average body temperature.

Sensory Modalities

Interrogating the average person regarding the human senses results in enumerating the five fundamental senses—gustation, olfaction, tactician, audition, and vision. Nevertheless, such a list is not exhaustive. Crucially absent is equilibrioception, or our sense of balance. Furthermore, the broad category of tactician can be categorized into more specialized modalities, such as pressure, vibration, stretch, and position of hair follicles, all discerned by different mechanoreceptors. Additional neglected modalities encompass thermoception, the detection of temperature by thermoreceptors, and nociception, the perception of pain via nociceptors.

In physiology, sensory perception can be categorized into general or specific frameworks. General sensory perception is pervasive throughout the body, with receptor cells embedded within the structure of other organs. For instance, mechanoreceptors in the skin, muscles, or blood vessel walls exemplify this category. Such general senses typically contribute to touch, proprioception (body's spatial orientation), kinesthesia (motion awareness), or autonomic functions through visceral senses. In contrast, specialized senses are associated with specific organs like the eye, inner ear, tongue, or nose.

Every sense is designated as a sensory modality, a term that encapsulates the concept of information encoding and mirrors the idea of transduction. The primary sensory modalities can be cataloged based on their respective transduction mechanisms. Taste and smell fall under chemical senses, while touch, often denoted as a general sense, encompasses chemical sensation in the form of nociception or pain. Mechanoreceptors detect sensations like pressure, vibration, muscle stretch, and hair movement due to an external stimulus. Mechanoreceptors also enable hearing and balance, while photoreceptors facilitate vision.

The myriad of sensory modalities, potentially amounting to 17, necessitates the fragmentation of the five primary senses into more specific subcategories or submodalities. Each sensory modality corresponds to the sensation of a unique stimulus. For instance, somatosensation, the general sense of touch, can be subdivided into various submodalities such as light pressure, deep pressure, vibration, itch, pain, temperature, or hair movement.

 Core: Anatomy and Physiology

Secondary Messengers in Hormone Action

JoVE 14976

Water-soluble hormones cannot cross the plasma membrane, so they rely on protein receptors that span the membrane to trigger intracellular signaling pathways. These pathways then activate second messengers inside the cell, including cAMP or calcium ions.

Many hormones bind to transmembrane G protein-coupled receptors that connect to regulatory G proteins. These G proteins can then activate enzymes such as adenylyl cyclase or phospholipase C. Adenylyl cyclase converts ATP to cAMP, activating protein kinases. Kinases are enzymes that add phosphate groups to other proteins, initiating a phosphorylation cascade. Epinephrine, glucagon, parathyroid hormone, and luteinizing hormone all act through cAMP. For instance, epinephrine activates the β-adrenergic receptor, initiating the GPCR signaling cascade to regulate the "fight or flight" response.

Phospholipase C, conversely, breaks down the membrane phospholipid PIP2 into IP3 and DAG. IP3 then moves to the endoplasmic membrane and attaches to an IP3-gated calcium channel, which causes a release of calcium ions into the cytoplasm. The change in calcium concentration can induce physiological effects such as muscle contraction. Oxytocin and hypothalamic regulatory hormones amplify signals through PIP2 and calcium ions.

 Core: Anatomy and Physiology

Hormones Regulating Blood Glucose

JoVE 14992

Insulin is released by beta cells of the pancreas when blood glucose levels are high. It facilitates glucose absorption and utilization in insulin-dependent cells with insulin receptors on their plasma membranes. Insulin promotes glucose uptake by increasing the number of glucose transport proteins in the cell membrane, allowing glucose to enter the cell. As a result, glucose utilization and ATP production are enhanced.

In addition to accelerating glucose uptake and utilization, insulin has other effects. It stimulates glycogen formation, a glucose storage polysaccharide, in skeletal muscle fibers and liver cells for later use. Insulin also promotes amino acid absorption and protein synthesis in all target cells, preventing the conversion of amino acids into glucose. Furthermore, insulin stimulates the formation of triglycerides in adipocytes, facilitating the absorption of fatty acids and glycerol.

However, it is important to note that not all cells in the body are insulin-dependent. Cells in the brain, kidneys, lining of the digestive tract, and red blood cells lack insulin receptors but can still absorb and use glucose independently of insulin stimulation.

When blood glucose levels drop below normal, the pancreas releases glucagon from alpha cells to mobilize energy reserves. Glucagon binds to its receptor on the target cell's plasma membrane, activating adenylate cyclase and producing cAMP as a second messenger. The primary effects of glucagon include stimulating the breakdown of glycogen in skeletal muscle fibers and liver cells releasing glucose molecules for energy metabolism or into the bloodstream. Glucagon also stimulates the breakdown of triglycerides in adipocytes, releasing fatty acids for use by other tissues. Additionally, it stimulates glucose production and release from liver cells through gluconeogenesis, contributing to an increase in blood glucose concentration.

The secretion of glucagon and insulin is regulated by pancreatic alpha and beta cells, respectively, in response to changes in blood glucose levels.

 Core: Anatomy and Physiology

Norton's Theorem

JoVE 15054

Norton's theorem is a fundamental principle stating that a linear two-terminal circuit can be substituted with an equivalent circuit, which comprises a current source (N) in parallel with a resistor (RN). Here, N represents the short-circuit current flowing through the terminals, and RN stands for the input or equivalent resistance at the terminals when all independent sources are deactivated. This implies that the circuit illustrated in Figure (a) can be exchanged with the one depicted in Figure (b).

Figure1

Figure (a)

Figure1

Figure (b)

The determination of Norton's resistance involves the process of setting all independent sources within the circuit to zero. The resulting value represents Norton's resistance, which is essential for the theorem's application. To calculate Norton's current, one must restore all sources to their original configurations and then ascertain the short-circuit current between the marked terminals. Replacing the network with the Norton equivalent circuit and reconnecting the load resistor is necessary to complete the transformation. To establish either the Thévenin or Norton equivalent circuit, one must have knowledge of the open-circuit voltage, short-circuit current, and input or equivalent resistance. These parameters serve as crucial inputs for these circuit transformations. The close relationship between Norton's and Thévenin's theorems given by the following relations,

Equation1

Equation1

coupled with the principles of Ohm's law, provides valuable tools for resolving intricate electrical circuits.

 Core: Electrical Engineering

Characteristics of OpAmp

JoVE 15077

The operational amplifier, commonly known as an op-amp, is a specially designed electronic circuit component. Its purpose is to work in conjunction with other circuit elements to execute a defined signal-processing operation. Consider an equivalent circuit model of an op-amp, as depicted in Figure 1; the output section comprises a voltage-controlled source in parallel with the output resistance Ro.

Figure1

Figure 1: The equivalent circuit of the nonideal op-amp

Figure 1 shows that the input resistance Ri equates to the Thėvenin equivalent resistance observable at the input terminals. Similarly, the output resistance Ro corresponds to the Thėvenin equivalent resistance discernible at the output. The differential input voltage (vd) is calculated as the difference between the noninverting terminal's voltage relative to the ground and the inverting terminal's voltage relative to the ground.

The op-amp discerns the difference between these two inputs and amplifies it by the gain A, thus generating a voltage that appears at the output. Here, A represents the open-loop voltage gain - named so because it signifies the gain of the op-amp in the absence of any external feedback from the output to the input.

Feedback is an integral concept to comprehend when studying op-amp circuits. Negative feedback can be achieved by feeding the output back to the op-amp's inverting terminal. When a feedback path from output to input exists, the resulting ratio of the output voltage to the input voltage is known as the closed-loop gain.

The op-amp does have its practical limitations. One such limitation is that the magnitude of its output voltage cannot surpass |VCC|. Essentially, the power supply voltage determines and restricts the output voltage. Depending on the differential input voltage, the op-amp can function in one of three modes: positive saturation, linear region, or negative saturation. If an attempt is made to increase vd beyond the linear range, the op-amp reaches saturation, yielding either vo = VCC or vo = -VCC.

 Core: Electrical Engineering

RC Circuit with Source

JoVE 15094

When a DC source is abruptly applied to an RC (Resistor-Capacitor) circuit, the voltage can be represented as a unit step function. The voltage across the capacitor, known as the step response, characterizes how the circuit reacts to this sudden change in input.

Due to the inherent properties of a capacitor, its voltage cannot change instantaneously. This means that immediately after the switch is closed, the capacitor's voltage remains the same as it was just before the switch was closed.

By applying Kirchhoff's current law at the moment the switch is closed (t=0), rearranging the terms, and rewriting the equation for times t>0, a first-order differential equation is obtained. This equation describes how the current through the circuit changes with time following the abrupt application of the DC source.

This differential equation is solved by integrating it, applying the limits, and taking the exponential on both sides. This yields the step response of the capacitor for times t>0. Combining this step response with the initial voltage across the capacitor (for t<0) gives the complete response of the RC circuit.

As time progresses, the voltage across the capacitor increases exponentially and approaches the applied source voltage. This process represents the charging of the capacitor.

If the capacitor is initially uncharged, the complete response of the circuit gets modified accordingly. From this modified response, the current through the capacitor is calculated. This current is observed to decrease exponentially with time, representing the gradual charging of the capacitor until it reaches the source voltage.

In conclusion, understanding the step response of an RC circuit provides valuable insights into how these circuits respond to sudden changes in input voltage. This knowledge is essential for designing and analyzing electronic circuits, particularly in applications such as signal processing, where capacitors are used extensively to filter or shape signals.

 Core: Electrical Engineering

Mesh Analysis for AC Circuits

JoVE 15111

In the domain of radio communication, the significance of impedance matching must be considered. It is crucial to ensure the efficient transmission of signals between radio transmitters and receivers. Achieving this balance involves using impedance-matching circuits, with one fundamental configuration comprising a resistor, capacitor, and inductor.

The process of harmonizing these impedances begins with a clear understanding of the input and output signals. Once these signals are known, the next step is calculating the current flowing through the capacitor in this circuit.

The angular frequency, extracted from the time-domain expression of the input voltage, assumes a critical role. It is a guiding factor in determining the impedance values of the inductor and the capacitor.

Equation1

Equation2

The circuit is then transformed into the frequency domain. This representation includes impedances, input and output signals, all expressed in polar form, simplifying the analysis. To delve deeper into the circuit's operation, mesh currents are assigned, and Kirchhoff's voltage law (KVL), a foundational principle in mesh analysis, is applied. Importantly, mesh analysis is particularly suited for planar circuits.

The outcome of this meticulous analysis yields a set of linear simultaneous equations, which can be elegantly represented in matrix form. Cramer's rule comes into play to reveal the mesh currents, allowing for the determination of the current shared across the capacitor.

Substituting the calculated mesh currents provides the current flowing through the capacitor, initially expressed in polar form. As a result, this data is skillfully transformed into the time domain, understanding and optimizing the impedance-matching circuit.

 Core: Electrical Engineering

Acid–Base Titration: Overview

JoVE 17359

An acid-base titration is a technique used to determine the concentration of an unknown acid or base, using a titrant of known concentration–either a base for acid titration or an acid for base titration. The process involves gradually adding the titrant, leading to a predictable change in the pH of the solution. This change is plotted on a titration curve, showing how a solution's pH varies with the amount of titrant added. Such curves are instrumental in monitoring the titration's progress and identifying the endpoint.

Titration curves can be formed in various combinations: between strong acids and strong bases, strong acids and weak bases, or weak acids and strong bases. These curves typically display an S-shape, but their specific characteristics vary depending on the strength of the acid or base involved. For strong acids neutralized by strong bases, the equivalence point - where the moles of acid equal the moles of base - is reached at a pH of 7, indicating complete neutralization. In the case of weak acids, this equivalence point occurs above pH 7, reflecting the transformation of the weak acid into its conjugate base. The pH changes more significantly near the equivalence point in strong acid titrations compared to weak acid titrations. In the latter, there is also a buffer region and a half-equivalence point, at which the pH equals the pKa of the acid. The concentrations of the weak acid and its conjugate base are equal—on the other hand, titrating strong or weak bases with strong acid results in an inverted S-shaped curve, starting from a higher pH and decreasing with the addition of the strong acid.

 Core: Analytical Chemistry

Drosophila Development and Reproduction

JoVE 5093

One of the many reasons that make Drosophila an extremely valuable organism is that the molecular, cellular, and genetic foundations of development are highly conserved between flies and higher eukaryotes such as humans. Drosophila progress through several developmental stages in a process known as the life cycle and each stage provides a unique platform for developmental research. This video introduces each stage of the Drosophila life cycle and details the physical characteristics and major developmental events that occur during each stage. Next, the video discusses the genetic regulation of pattern formation, which is important for establishing the body plan of the organism and specifying individual tissues and organs. In addition, this video gives an overview of Drosophila reproduction, and how to use the reproductive characteristics of Drosophila to set up a genetic cross. Finally, we discuss examples of how the principles of Drosophila development and reproduction can be applied to research. These applications include RNA interference, behavioral assays of mating behaviors, and live imaging techniques that allow us to visualize development as a dynamic process. Overall, this video highlights the importance of understanding development and reproduction in Drosophila, and how this knowledge can be used to understand development in other organisms.

 Biology I

Development of the Chick

JoVE 5155

The chicken embryo (Gallus gallus domesticus) provides an economical and accessible model for developmental biology research. Chicks develop rapidly and are amenable to genetic and physiological manipulations, allowing researchers to investigate developmental pathways down to the cell and molecular levels.

This video review of chick development begins by describing the process of egg fertilization and formation within the chicken reproductive tract. Next, the most commonly used chick staging nomenclature, the Hamburger Hamilton staging series, is introduced. Major events in chick development are then outlined, including the dramatic cellular movements known as gastrulation that form the three major cell layers: The ectoderm, mesoderm, and endoderm. Cells from these layers go on to generate all the tissues within the organism, as well as extraembryonic membranes, which are necessary for the transport of gases, nutrients, and wastes within the eggshell. To conclude the discussion, some exciting techniques will be presented as strategies for studying chick development in greater detail.

 Biology II

Murine In Utero Electroporation

JoVE 5208

In utero electroporation is an important technique for studying the molecular mechanisms that guide the proliferation, differentiation, migration, and maturation of cells during neural development. Electroporation enables the rapid and targeted delivery of material into cells by utilizing electrical pulses to create transient pores in cell membranes. Although electroporation has traditionally been used in in vitro studies, scientific advancements have now broadened its utilization to intact organs, such as those found in mouse embryos developing in utero.

This video will introduce the key principles behind in utero electroporation in addition to reviewing the basic surgical techniques required to access developing embryos within a pregnant rodent. Details of the injection and electroporation steps are provided along with important considerations for directing gene delivery to specific brain regions. Finally, neurobiological applications of in utero electroporation are presented, such as investigating how specific genes contribute to neural development and how connections form between developing neurons.

 Neuroscience

Whole-Mount In Situ Hybridization

JoVE 5330

Whole-mount in situ hybridization (WMISH) is a common technique used for visualizing the location of expressed RNAs in embryos. In this process, synthetically produced RNA probes are first complementarily bound, or "hybridized," to the transcripts of target genes. Immunohistochemistry or fluorescence is then used to detect these RNA hybrids, revealing spatial and temporal patterns of gene expression. Unlike traditional in situ hybridization techniques, which require thin tissue sections whose images will need to be computationally reassembled, the whole-mount technique allows gene expression patterns to be assessed over the entire embryo or structure.

This video will introduce the basic concepts of whole mount staining and detail key procedural steps, including probe design and production, embryo fixation and staining, and post-hybridization signal detection. Viewers will then learn about how developmental biologists are applying WMISH to current research studies.

 Developmental Biology

Electro-encephalography (EEG)

JoVE 5420

EEG is a non-invasive technique that can measure brain activity. The neural activity generates electrical signals that are recorded by EEG electrodes placed on the scalp. When an individual is engaged in performing a cognitive task, brain activity changes and these changes can be recorded on the EEG graph. Therefore, it is a powerful tool for cognitive scientist aiming to better understand the neural correlates associated with different aspects of cognition, which will ultimately help them devise improved treatments for patients with cognitive deficits.

Here, JoVE presents a brief overview of EEG and its applications in cognitive research. First, we discuss where and how EEG signals are generated. Then, we explain the use of EEG in studying cognition along with a detailed step-by-step protocol to perform an EEG experiment. Lastly, the video reviews some specific cognitive experiments that use EEG in combination with other techniques such as functional Magnetic Resonance imaging (fMRI) or transcranial direct current stimulation (tDCS).

 Behavioral Science

An Overview of Genetic Engineering

JoVE 5552

Genetic engineering – the process of purposefully altering an organism’s DNA – has been used to create powerful research tools and model organisms, and has also seen many agricultural applications. However, in order to engineer traits to tackle complex agricultural problems such as stress tolerance, or to realize the promise of gene therapy for treating human diseases, further advances in the field are still needed. Important considerations include the safe and efficient delivery of genetic constructs into cells or organisms, and the establishment of the desired modification in an organism’s genome with the least “off-target” effects.

JoVE’s Overview of Genetic Engineering will present a history of the field, highlighting the discoveries that confirmed DNA as the genetic material and led to the development of tools to modify DNA. Key questions that must be answered in order to improve the process of genetic engineering will then be introduced, along with various tools used by genetic engineers. Finally, we will survey several applications demonstrating the types of experimental questions and strategies in the field today.

 Genetics

Cell-surface Biotinylation Assay

JoVE 5647

A cell can regulate the amount of particular proteins on its cell membrane through endocytosis, following which cell surface proteins are effectively sequestered in the cytoplasm. Once within a cell, these surface proteins can be either destroyed or “recycled” back to the membrane. The cell surface biotinylation assay provides researchers with a way to study these phenomena. The technique makes use of a derivative of the small molecule biotin, which can label surface proteins and then be chemically cleaved. However, if the surface protein is endocytosed, the biotin derivative will be protected from cleavage. Thus, by analyzing the uncleaved, endocytosed biotin label, scientists can assess the amounts of internalized surface proteins.

In this video, we review the concepts behind the biotinylation assay, delving into the chemical structure of the biotin derivative and the mechanism of its cleavage. This is followed by a generalized protocol of the technique, and finally, a description of how researchers are currently using it to study the dynamics of different cell surface proteins.

 Cell Biology

Protein Crystallization

JoVE 5689

Protein crystallization, obtaining a solid lattice of biomolecules, elucidates protein structure and enables the study of protein function. Crystallization involves drying purified protein under a combination of many factors, including pH, temperature, ionic strength, and protein concentration. Once crystals are obtained, the protein structure can be elucidated by x-ray diffraction and computation of an electron density model.

This video introduces protein crystallization and shows a general procedure. Protein expression and purification, crystallization, and x-ray diffraction are covered in the procedure. Applications of protein crystallization include in silico drug design, binding site determination, and membrane protein structure analysis.

Protein crystallization is the process of obtaining a latticed solid form of a protein. These crystals are especially valuable to structural biologists, assisting in the study of protein function. Other techniques, such as mass spec or SDS-PAGE, can only provide information on the one-dimensional structure of proteins. Protein crystallization is complemented by the techniques of recombinant protein expression and x-ray diffraction. This video will show the principles of protein crystallization, a general laboratory procedure, and several of its applications in the biochemical field.

The first step required in the process is to obtain milligram quantities of very pure protein, typically using recombinant protein expression. The gene corresponding to the protein of interest is cloned into an expression vector, and the expressed protein is fused to an affinity tag, such as poly-histidine, to assist in the purification by affinity chromatography. To learn more, see this collection's video on affinity chromatography.

Formation of the purified protein into crystals is dependent on the proper combination of many factors, including pH, ionic strength, concentrations of precipitant and protein, temperature, and rate of equilibration. The most common method used is vapor diffusion, of which there are two categories: hanging drop and sitting drop. A droplet containing pure protein, buffer, and precipitant, which is an ionic solid that binds water molecules, reducing water availability for the protein and mimicking higher protein concentration, is in an enclosed microwell with a reservoir with a more highly concentrated mixture of the same buffer and precipitant. At the beginning, the concentrations of protein and precipitant are too low to cause crystallization. During the course of the experiment, water vaporizes from the droplet and collects in the reservoir; a decrease in the amount of water in the droplet causes the system to become supersaturated, and nucleation, followed by crystallization, can occur. The net transfer of water from the droplet is in equilibrium, and the system is maintained until the process is complete.

To visualize the 3D structure, x-ray diffraction is used. To obtain x-ray data from a crystal, it is placed in a monochromatic x-ray beam, where it is exposed to the beam at all angles. Each exposure provides an image, where each spot is a diffracted x-ray, which emerges from the crystal and is registered by a detector. The data are combined to produce a model of the arrangement of atoms within the crystal. The resulting crystal structure demonstrates the 3-dimensional placement of the atoms, with a typical resolution of 2 angstroms.

Now that we have covered the principles of protein crystallization, let us look at a generalized protocol.

To begin the procedure an expression vector containing the gene of interest is transformed into cells. The cells are incubated and at mid-log phase, expression is initiated by adding an inducer, such as IPTG, which triggers transcription of the gene's mRNA. After protein expression, the crude material is suspended in lysis buffer, and then clarified by centrifugation.

The clarified lysate is then loaded onto a nickel column, and the polyhistidine-tagged protein binds to the column while all other biomolecules are washed away.

Once several milligrams of pure protein have been obtained, it is ready for crystallization by vapor diffusion. A 24-well hanging/sitting drop tray is filled with varying concentrations of sodium chloride and sodium acetate buffer solutions. For the sitting drop method, equal volumes of protein and reservoir solution are pipetted onto the shelf above each well, and then the tray is covered with transparent tape. The tray is then placed in an incubation chamber, and the wells are monitored for growth the following day, then every few days.

Once a proper crystal has been obtained it is ready for x-ray diffraction analysis. The crystal is mounted on a goniometer to position the crystal at selected orientations. The crystal is illuminated with a monochromatic beam of x-rays at all angles, producing a diffraction pattern. The software converts the two-dimensional images, taken at different orientations, to a three-dimensional model of the density of electrons within the crystal by determining the positions of the atoms in the crystal.

Now that we have reviewed a procedure, let's review some useful applications of protein crystallization, and another crystallization technique.

Protein crystallization may be used for in silico drug design. The three-dimensional structure of Influenza virus's polymerase basic protein 2, which has been linked to viral infection in mammals, was determined by crystallization and x-ray diffraction. Potential binding sites in the protein are visualized, and with the use of a docking program, a three-dimensional molecule was designed that would insert into a cleft in the protein.

Co-crystallization of protein-DNA complexes is also a useful technique. DNA-binding proteins modulate a wide variety of biological functions such as transcription and DNA polymerization and DNA repair; and crystal structures of these complexes can provide insight into protein function, mechanism, and the nature of the specific interaction. The E. coli protein SeqA, a negative regulator of DNA replication, was co-crystallized with hemimethylated DNA.

Integral membrane proteins such as G-protein coupled receptors, or GCPRs, are difficult to crystallize due to their limited amount of polar surface area available for forming crystal lattice contacts, which has led to the development of fusion-protein-assisted protein crystallization. Genes encoding β2 adrenergic receptor, a GCPR, and a lysozyme were inserted into an expression vector. The crystallization of the β2AR- lysozyme fusion protein was achieved due to the increased extracellular hydrophilic surface over the naturally hydrophobic β2AR, provided by the lysozyme, necessary for forming packing interactions in the crystal lattice.

You've just watched JoVE's video on protein crystallization. This video described its principles, a generalized protocol, and some its uses in the biomedical field. Thanks for watching!

 Biochemistry

Microfabrication via Photolithography

JoVE 5789

The fabrication of BioMEM devices is often done using a microfabrication technique called photolithography. This widely used method utilizes light to transfer a pattern onto a silicon wafer, and provides the basis for the fabrication of many types of BioMEM devices.

This video presents the photolithography technique, shows how the process is performed in the cleanroom, and introduces some applications of the process.

 Bioengineering

Effect of Temperature Change on Reaction Rate

JoVE 11698

The Arrhenius equation,

Figure1

relates the activation energy and the rate constant, k, for many chemical reactions.

In this equation, R is the ideal gas constant, which has a value 8.314 J/mol·K, T is the temperature in kelvin, Ea is the activation energy in joules per mole, e is the constant 2.7183, and A is a constant called the frequency factor, which is related to the frequency of collisions and the orientation of the reacting molecules.

The frequency factor, A, reflects how well the reaction conditions favor correctly oriented collisions between reactant molecules. An increased probability of effectively oriented collisions results in larger values for A and faster reaction rates.

The exponential term, e−Eₐ/RT, describes the effect of activation energy on the reaction rate. According to kinetic molecular theory, the temperature of matter is a measure of the average kinetic energy of its constituent atoms or molecules—a lower activation energy results in a more significant fraction of adequately energized molecules and a faster reaction.

The exponential term also describes the effect of temperature on the reaction rate. A higher temperature represents a correspondingly greater fraction of molecules possessing sufficient energy (RT) to overcome the activation barrier (Ea). This yields a higher value for the rate constant and a correspondingly faster reaction rate.

The minimum energy necessary to form a product during a collision between reactants is called the activation energy (Ea). The difference in activation energy required and the kinetic energy provided by colliding reactant molecules is a primary factor affecting the rate of a chemical reaction. If the activation energy is much larger than the average kinetic energy of molecules, the reaction will occur slowly, since only a few fast-moving molecules will have enough energy to react. If the activation energy is much smaller than the molecules' average kinetic energy, a large fraction of molecules will be adequately energetic, and the reaction will proceed rapidly.

Reaction diagrams are widely used in chemical kinetics to illustrate various properties of a reaction of interest. It shows how a chemical system's energy changes as it undergoes a reaction, converting reactants to products.

This text is adapted from Openstax, Chemistry 2e, Section 12.5: Collision Theory.

 Core: Organic Chemistry

Solvating Effects

JoVE 11739

An understanding of the solvating effect helps rationalize the relation between solvation and acidity of the compound. In addition, this also explains the relative stability of conjugate bases for compounds with different pKa values. This lesson details, in-depth, the principle of solvating effects. The strength of an acid and the stability of its corresponding conjugate base are determined using pKa values. This observed relationship is a consequence of solvation, which is the interaction between a dissolved ion and solvent molecules. During this process, the solvent molecules surround the ions and stabilize them.

Solvation of dissolved ions can be classified into three types: (i) donor interaction, (ii) charge–dipole interaction, and (iii) hydrogen-bonding interaction. In the donor interaction, a solvent donates its unshared electron pairs to the dissolved ion. The solvent acts as a Lewis base, and the ion acts as a Lewis acid. In the second type, charge–dipole interactions are observed in polar solvents, where their dipole moments can interact with the charged ions. This involves rearranging the positive partial charge on the solvent molecules to align with the negative charge of the ions, thus stabilizing the ions. For instance, as noted in the solvation of ethanol, the ethoxide anion, which is the conjugate base, is solvated by the positive center of the solvent’s dipole that stabilizes it effectively. Lastly, when the ions are stabilized by hydrogen bonding between the solvent molecules and the dissolved ions, the interaction is called a hydrogen-bonding interaction.

The interactions between the dissolved ions and solvent molecules influence their stability, which is directly proportional to the strength of the acidity.  Accordingly, the stability of such ions increases with a larger number of interactions when they are surrounded by more solvent molecules. Therefore, during solvation, the steric hindrance from bulky substituents on the molecule plays an important role. Compounds with less bulky groups are sterically unhindered, allowing for more interaction with solvent molecules.

In contrast, the compounds that possess bulky groups have steric hindrance and are consequently poorly solvated. As a result, the sterically unhindered ion demonstrates more stability, making its corresponding acid stronger. This is demonstrated with the comparison of acidity of ethanol, isopropanol, and tert-butanol. With the increasing size of substituents, the corresponding conjugate base of each of these compounds has more steric hindrance. Hence, it is less solvated. As a result, isopropanol is a weaker acid (pKa=17.10) than ethanol (pKa=16.00), and tert-butanol (pKa=19.20) is a weaker acid than isopropanol (pKa=17.10). In summation, the steric hindrance of the conjugate base anions defines the degree of solvation. Low solvation leads to instability of the dissolved ion that makes the corresponding acid weak. 

 Core: Organic Chemistry

Acid-Catalyzed Hydration of Alkenes

JoVE 11776

Alkenes react with water in the presence of an acid to form an alcohol. In the absence of acid, hydration of alkenes does not occur at a significant rate, and the acid is not consumed in the reaction. Therefore, alkene hydration is an acid-catalyzed reaction.

Figure1

Strong acids, such as sulfuric acid, dissociate completely in an aqueous solution, and the acid participating in the reaction is the hydronium ion.

Figure2

The first step is the slow protonation of an alkene at the less-substituted end to form the more-substituted carbocation.

Figure3

The second step is the nucleophilic attack by water at the carbocation to give an oxonium ion.

Figure4

In the last step, water, with a pKa of 15.7, acts as a base and deprotonates the acidic oxonium ion (protonated alcohol), which has a pKa of approximately –2, to yield the final product.

Figure5

The two processes, hydration of alkenes to form alcohols and the dehydration of alcohols to form alkenes, are in equilibrium with each other. The control over this equilibrium can be explained by Le Chatelier’s principle, which states that a system at equilibrium will adjust to minimize any stress placed on the system.

In the hydration of 2-methylpropene, water is on the left side of the reaction. When the amount of water increases, the equilibrium shifts towards the right, producing more alcohol. In contrast, eliminating water from the system shifts the equilibrium to produce more alkene. Thus, the presence of dilute acids favors the formation of alcohols from alkenes, while the reverse occurs in the presence of concentrated acids that contain very little water.

Addition reactions are temperature-dependent. The enthalpy term for these reactions is negative as new bonds are formed during the process. In contrast, the entropy term is positive as the two reactant molecules give one molecule of product.

At low temperatures, the entropy term is small and the enthalpy term dominates. Thus, the Gibbs free energy is negative, and the equilibrium constant being greater than one promotes the formation of product over reactants.

Figure7

However, at high temperatures, the large entropy term dominates the enthalpy term and the Gibbs free energy is positive. The equilibrium constant being less than one reverses the reaction, implying that reactants will be favored over products.

Figure7

 Core: Organic Chemistry

The Role of Actin and Myosin in Non-muscle Cells

JoVE 11804

Actin and myosin or actomyosin filaments also play a significant role in cells other than those involved in muscle contraction (which occurs within the sarcomere of muscle cells). The mechanism of non-muscle cell contractile bundles was first observed in Dictyostelium and Acanthamoeba. In non-muscle cells, two bundles are commonly found: stress fibers and actomyosin adherence belts. These contractile bundles are smaller and less organized than the ones found in muscle cells. They  are held together by accessory proteins such as  fascin, filamin, and fimbrin depending on where they are located. The formation and contractile ability of these actomyosin bundles are regulated by phosphorylation.

In non-muscle cells,  actin and non-muscle myosin II (NMM) filaments combine to form stress fibers. There are about 15 to 20 myosin filaments in each bundle. The mechanism of action of these actinomyosin contractile bundles is similar to those in muscle fibers, where  ATP hydrolysis by the myosin globular head drives the actin contraction. The stress fibers, along with the focal adhesions present at the cell edges, regulate the mechanosensitive machinery in the cells. These bundles are also responsible for regulating cell morphogenesis. In cells, these stress fibers are found near the cell edges, where they help the cell anchor on the substratum.

In epithelial cells, the actomyosin bundles are attached to the adhesion junctions as circumferential or adherence belts close to the plasma membrane. These belts regulate the apical constriction of polarized  epithelial cells. In dividing cells, these actomyosin bundles are recruited by the septin ring at the cell cleavage furrow to form the contractile ring, which is a major step in  cytokinesis.

 Core: Cell Biology

The Movement of Organelles and Vesicles

JoVE 11912

In eukaryotic cells,  cytoskeletal filaments such as actin, microtubules, and intermediate filaments form a mesh-like cytoskeletal network. These filaments serve as tracks for transporting cellular cargo. Specialized motor proteins use the chemical energy stored in adenosine triphosphate (ATP) for this transport. During interphase, microtubules are polarized, with the plus-end towards the cell periphery and the minus-end towards the cell center. Two microtubule-associated motor proteins, kinesin and cytoplasmic dynein, transport organelles and vesicles.

Transport via kinesins

Kinesin is a plus-end-directed, microtubule-associated motor protein made up of two heavy chains and two light chains. Kinesin motors are highly efficient as they undergo hundreds of ATP hydrolysis cycles without dissociating from the microtubules while transporting cellular cargos. During interphase, kinesins transport organelles and vesicles across the cell. The kinesins attach to specific vesicles or organelles with the help of the receptor domain present in its light chain. Some kinesin motor proteins have also shown preferential binding for post-translationally modified microtubules. For example,  acetylation and detyrosination of microtubules alter the binding and mobility of kinesin-1. This selective binding plays a crucial role in the axonal cargo transport within the neurons.

Transport via dyneins

Dyneins are minus-end-directed motor proteins. These motors are responsible for carrying out crucial functions within the cell, such as the arrangement of axonal complex microtubule assemblies, signal transduction within the primary cilium, and positioning of Golgi apparatus near the cell center. Among the two dyneins present in the cell, cytoplasmic dynein is primarily responsible for transporting organelles and vesicles. Dynein has two globular heads attached to two stalks and intermediate and light chains that form the upper domain. ATP hydrolysis takes place within the globular heads. The movement of cytoplasmic dynein resembles walking, where one stalk moves ahead, followed by the other. The large cytoplasmic dynein molecule cannot bind directly to the organelles or vesicles; it requires an additional protein called dynactin. Dynactin has a short actin-like filament, Arp1. The Arp-1 domain recognizes and attaches to the receptor domain of the organelles or vesicles. The cytoplasmic dynein and dynactin are activated only after binding with a specific organelle or vesicle, forming a tripartite complex. The tripartite complex then transports the organelles and vesicles.

 Core: Cell Biology

Sharpless Epoxidation

JoVE 12108

The conversion of allylic alcohols into epoxides using the chiral catalyst was discovered by K. Barry Sharpless and is known as Sharpless epoxidation. The use of a chiral catalyst enables the formation of one enantiomer of the product in excess. This chiral catalyst is mainly a chiral complex of titanium tetraisopropoxide and tartrate ester (specific stereoisomer). The stereoisomer used in the chiral catalyst dictates the formation of the enantiomer of the product. In other words, the use of L-(+)-diethyl tartrate leads to enantiomers having the epoxide ring below the plane, while with D-(−)-diethyl tartrate, to enantiomers with the epoxide ring above the plane. The high enantioselectivity of the reaction can be explained by considering the activation energies required for the reaction to proceed in the forward direction in the presence of the chiral catalyst. As shown in Figure 1, compared to the uncatalyzed reaction (blue curve), the activation energy of the reaction decreases dramatically with the addition of the chiral catalyst (red and green curves). Moreover, the activation energy for the formation of one enantiomer (red curve) is lowered more than that of another enantiomer (green curve), leading to the formation of one enantiomer in excess. Hence, Sharpless epoxidation reaction can be utilized for the synthesis of desired enantiomers of the product.

Figure1

The stereochemistry of the product formed when any allylic alcohol is subjected to Sharpless epoxidation can be predicted by simply orienting the allylic alcohol molecule in a plane with the hydroxyl groups pointing towards the lower right corner, as shown in Figure 2. On this planar structure, D-(−)-diethyl tartrate delivers the oxygen from the top face of the alkene, making the epoxide formation feasible from above the plane, while L-(+)-diethyl tartrate delivers the oxygen from the bottom face of the alkene, thereby installing the epoxide ring from below the plane.

Figure2

 Core: Organic Chemistry

Cytoskeletal Coordination in Cell Migration

JoVE 12260

A migrating cell changes its shape during the cyclic events of attachment and detachment from the substratum and repositions the cell organelles correspondingly. These complex events are orchestrated by the dynamic cytoskeletal network comprising actin filaments, intermediate filaments, and microtubules. Cytoskeletal crosstalk — the direct and indirect communication between the different components — is crucial for this coordination. Direct communication involves various linker proteins that form cross-bridges between the components. For example, the linker protein spectraplakin bridges microtubules with actin filaments, and KASH proteins link the nucleus to the cytoskeleton. Indirect communication is mediated via signaling cascades, such as those involving the Rho proteins. 

Rho, the Master Regulators

The small GTPases of the Rho family proteins, such as RhoA, Cdc42, and Rac1, are the master regulators in establishing cell polarity by acting on all three cytoskeletal components. For example, Cdc42 not only directs actin reorganization at the leading edge but also reorganizes intermediate filaments by regulating their transport on the microtubules. The microtubules also direct the transport of various proteins and vesicles, which, in turn, regulate actin dynamics at the leading edge. Such positive feedback helps maintain the synchronized polarity of the different cytoskeletal components.

 Core: Cell Biology

Mechanism of Angiogenesis

JoVE 12509

Blood vessel formation starts early during embryonic development, around day 7. In the extraembryonic yolk sac, mesodermal precursor cells called hemangioblast proliferate and differentiate into angioblast. Angioblasts express vascular endothelial growth factor receptor 2 or VEGFR2, which binds VEGF-A, a proangiogenic factor, guiding blood vessel formation. VEGF signaling promotes angioblasts to form a blood island in the developing embryo. Angioblasts further differentiate, giving rise to endothelial cells, which aggregate and form the primitive vascular network called the vascular plexus.

    Following vasculogenesis, the vascular plexus or the preexisting blood vessels guide new blood vessel formation through angiogenesis. Angiogenesis occurs by two distinct mechanisms: sprouting angiogenesis and intussusceptive angiogenesis.

  1. Sprouting angiogenesis:

In sprouting angiogenesis, angiogenic stimuli such as injury or growth spurt induce VEGF secretion from surrounding endothelial cells. VEGF binds VEGFR on these cells and signals them to produce matrix metalloproteases, which helps in basement membrane degradation following the detachment of pericytes from the vessel wall. VEGF and other angiogenic signals also stimulate endothelial cells to differentiate into tip and stalk cells. The tip cells extend filopodial structures resembling blind end tubes called sprouts. Tip cells migrate towards the angiogenic stimulus following the VEGF gradient. The stalk cells behind the angiogenic tip proliferate, elongating the tubular structure. Next, the newly formed blood vessel develops a lumen by one of the two processes: cell hollowing and cord hollowing. In the cell hollowing method, intracellular vacuoles fuse and connect adjacent cells, forming a continuous lumen. In the cord hollowing method, endothelial cells forming the tube change their shape to develop a central tubular lumen towards their extracellular side.

  1.  Intussusceptive angiogenesis:

Intussusceptive angiogenesis involves splitting a vessel into two, also called splitting angiogenesis. It is a faster way of new blood vessel formation. The newly formed vessels join the existing vessels and complete the vascular network. Pericytes and smooth muscles surround the newly formed vessel stabilizing them as blood flows through them.

 Core: Cell Biology

Factors Affecting α-Alkylation of Ketones: Choice of Base

JoVE 13072

α-Alkylation of ketones is achieved in the presence of alkyl halides and a base. The reaction proceeds via the formation of an enolate ion followed by nucleophilic substitution. The choice of base employed is essential as it is the key factor in determining the reaction outcome.

The reaction involving bases like EtO whose conjugate acid EtOH (pKa = 15.9) is stronger than the ketone (pKa = 19.2) results in an equilibrium mixture with higher ketone concentration. As a consequence, side reactions become predominant over α-alkylation. Using bases like LDA, whose conjugate acid NH(CHMe2)2 is weaker (pKa = 36) than the ketones, leads to an irreversible enolate ion formation, excluding undesirable side reactions. Hence, the nucleophilic enolate further undergoes substitution with alkyl halides to produce the desired α-alkylated ketone. 

 Core: Organic Chemistry

Hybridoma Technology

JoVE 13374

Hybridoma technology is used for the large-scale production of monoclonal antibodies. Monoclonal antibodies bind to only a single antigenic determinant or epitope. Such antibodies are used in research, diagnostics, and disease therapy. The hybridoma technology established in 1975 by Georges Köhler and Cesar Milstein was awarded the Nobel Prize in Medicine in 1984 for revolutionizing research and therapy.

Hybridoma Selection

Commonly used fusion techniques — electroporation, polyethylene glycol (PEG) mediated fusion, and the use of fusogenic viruses — generally produce hybridomas at a very low frequency. Fusion products thus have a high ratio of self-fused B-cells and myeloma cells, or unfused cells, with a low number of hybrid cells.

The hypoxanthine-aminopterin-thymidine (HAT) medium is used to select hybridomas by selectively allowing their growth. The aminopterin blocks the default nucleotide synthesis in cells. However, cells can utilize hypoxanthine and thymidine from the medium to synthesize nucleotides via the salvage pathway.

Myeloma cells deficient in the HGPRT enzyme cannot synthesize nucleotides via this pathway and thus, do not grow in HAT medium. Although B-cells produce functional HGPRT enzymes, they cannot divide indefinitely. Thus, only the hybrid cells with functional HGPT from the B-cells and immortality from the myeloma line can grow on the HAT selection medium. Such a culture of hybrid cells is called a hybridoma, which can be screened for monoclonal antibody production.

Hybridoma Screening

Hybridoma cultures from the HAT selection are plated onto a 96-well plate; each well containing only one cell. Using Enzyme-linked Immunosorbent Assay (ELISA), each cell is screened for the production of antibodies specific to the epitope of interest. Cells that test positive are then selected and grown in a larger culture vessel to establish hybridoma cell lines. These cell lines are a permanent source of unlimited monoclonal antibodies. Thus hybridoma technology is a convenient and cost-effective method for mass production of monoclonal antibodies.

 Core: Cell Biology

Phase Contrast and Differential Interference Contrast Microscopy

JoVE 13390

Phase-Contrast Microscopes

In-phase-contrast microscopes, interference between light directly passing through a cell and light refracted by cellular components is used to create high-contrast, high-resolution images without staining. It is the oldest and simplest type of microscope that creates an image by altering the wavelengths of light rays passing through the specimen. Altered wavelength paths are created using an annular stop in the condenser. The annular stop produces a hollow cone of light focused on the specimen before reaching the objective lens. The objective contains a phase plate with a phase ring. As a result, light traveling directly from the illuminator passes through the phase ring while light refracted or reflected by the specimen passes through the plate. This causes waves traveling through the ring to be about one-half of a wavelength out of phase with those passing through the plate.

Because waves have peaks and troughs, they can add together (if they are in-phase) or cancel each other out (if out-of-phase). When the wavelengths are out-of-phase, wave troughs cancel out wave peaks, which is called destructive interference. Structures that refract light then appear dark against a bright background of only unrefracted light. More generally, structures that differ in features such as refractive index will differ in levels of darkness. As it increases contrast without requiring stains, phase-contrast microscopy is often used to observe live specimens. Specific cellular structures, such as organelles in eukaryotic cells and endospores in prokaryotic cells, are well visualized with phase-contrast microscopy.

Differential Interference Contrast Microscopes

Differential interference contrast (DIC) microscopes (also known as Nomarski optics) are similar to phase-contrast microscopes in that they use interference patterns to enhance the contrast between different features of a specimen. In a DIC microscope, two beams of light are created in which the direction of wave movement (polarization) differs. Once the beams pass through either the specimen or specimen-free space, they are recombined. The effects of the specimens cause differences in the interference patterns generated by combining the beams. This results in high-contrast images of living organisms with a three-dimensional appearance. These microscopes are especially useful in distinguishing structures within live, unstained specimens.

 Core: Cell Biology

What is the Cell Cycle?

JoVE 13406

The cell cycle refers to the sequence of events occurring throughout a typical cell’s life. In eukaryotic cells, the somatic cell cycle has two stages: the interphase and the mitotic phase. During interphase, the cell grows, performs its basic metabolic functions, copies its DNA, and prepares for mitotic cell division. Then, during mitosis and cytokinesis, the cell divides its nuclear and cytoplasmic materials, respectively. This generates two daughter cells that are identical to the original parent cell. The cell cycle is essential for the growth of the organism, replacement of damaged cells, and regeneration of aged cells. Cancer is the result of uncontrolled cell division sparked by a gene mutation.

Cell Cycle Checkpoints

There are three major checkpoints in the eukaryotic cell cycle. At each checkpoint, the progression to the next cell cycle stage can be halted until conditions are more favorable. The G1 checkpoint is the first of these, where a cell’s size, energy, nutrients, DNA quality, and other external factors are evaluated. If the cell is deemed inadequate, it does not continue to the S phase of interphase. The G2 checkpoint is the second checkpoint. Here, the cell ensures that all of the DNA has been replicated and is not damaged before entering mitosis. If any DNA damage is detected that cannot be repaired; the cell may undergo apoptosis, or programmed cell death. The M or spindle checkpoint ensures that all the sister chromatids are correctly attached to the spindle microtubules at the metaphase plate before the cell enters anaphase.

Cancer: When the Cell Cycle Goes Awry

Cell cycle checkpoints ensure that healthy cells move through the cell cycle in a regulated way. However, cancer cells often bypass these checkpoints. Each successive round of unchecked cell division produces more damaged daughter cells. Additionally, cancer cells in the human body can divide many more times than normal cells, which can only undergo about 40-60 rounds of division. Cancer cells express telomerase, an enzyme that repairs the wear and tear at the ends of chromosomes that is typically caused by cell division.

 Core: Cell Biology

Crossing Over

JoVE 13438

Crossing over is the exchange of genetic information between homologous chromosomes during prophase I of meiosis I. Genetic recombination gives rise to allelic diversity in the newly formed daughter cells. In humans, crossing over produces genetically distinct haploid egg and sperm cells that undergo fertilization to produce unique offspring. Before cell division starts, the germ cell’s chromosome(s) undergo duplication in the S phase of the cell cycle. As the cells enter prophase I, duplicated chromosomes condense and form two sister chromatids (identical copies of the original chromosome) joined by the centromere. Next, the homologous chromosomes pair up and align the same gene segment from the maternal and paternal chromosomes, forming a synapse. A protein complex called the synaptonemal complex is formed that holds these homologs together. As crossing over proceeds, random pieces of DNA are swapped between the homologs, producing new combinations of alleles via homologous recombination. The 'chiasmata' mark the areas where the crossover of genetic information has occurred. As the synaptonemal complex begins to dissolve, the chiasma holds the homologous chromosomes together until recombination is completed and chromosomes are segregated correctly into the daughter cells.

 Core: Cell Biology

Stem Cell Culture

JoVE 13470

Stem cell research aims to find ways to use stem cells to regenerate and repair cellular damage. Over time, most adult cells undergo the wear and tear of aging and lose their ability to divide and repair themselves. Stem cells do not display a particular morphology or function. Adult stem cells, which exist as a small subset of cells in most tissues, keep dividing and can differentiate into a number of specialized cells generally formed by that tissue. These cells enable the body to renew and repair body tissues.

The mechanisms that induce a non-differentiated cell to become a specialized cell are poorly understood. In a laboratory setting, it is possible to induce stem cells to differentiate into specialized cells by changing the physical and chemical conditions of growth. Several sources of stem cells are used experimentally and are classified according to their origin and potential for differentiation. Human embryonic stem cells (hESCs) are extracted from embryos and are pluripotent. The adult stem cells that are present in many organs and differentiated tissues, such as bone marrow and skin, are multipotent, being limited in differentiation to the types of cells found in those tissues.

The stem cells isolated from umbilical cord blood are also multipotent, as are cells from deciduous teeth (baby teeth). Researchers have recently developed induced pluripotent stem cells (iPSCs) from mouse and human adult stem cells. These cells are genetically reprogrammed multipotent adult cells that function like embryonic stem cells; they are capable of generating cells characteristic of all three germ layers.

Because of their capacity to divide and differentiate into specialized cells, stem cells offer a potential treatment for diseases such as diabetes and heart disease. Cell-based therapy refers to treatment in which stem cells induced to differentiate in a growth dish are injected into a patient to repair damaged or destroyed cells or tissues. Many obstacles must be overcome for the application of cell-based therapy. Although embryonic stem cells have a nearly unlimited range of differentiation potential, they are seen as foreign by the patient’s immune system and may trigger rejection. Also, the destruction of embryos to isolate embryonic stem cells raises considerable ethical and legal questions.

In contrast, adult stem cells isolated from a patient are not seen as foreign by the body, but they have a limited range of differentiation. Some individuals bank the cord blood or deciduous teeth of their child, storing away those sources of stem cells for future use, should their child need it. Induced pluripotent stem cells are considered a promising advance in the field because using them avoids the legal, ethical, and immunological pitfalls of embryonic stem cells.

This text is adapted from openstax Anatomy and physiology 2e, Section 3.6: Cell differentiation.

 Core: Cell Biology

Methods of Nuclear Reprogramming

JoVE 13559

Nuclear reprogramming is a process of transforming one cell type into an unrelated cell type by epigenetic changes that alter the cell’s original gene expression pattern. Such epigenetic changes force cells to express a different set of genes, which play a significant role in inducing transformation into other cell types. Nuclear reprogramming offers applications in reproductive cloning for livestock propagation and regenerative medicine — developing patient-specific cells for injury repair.

Epigenetic Changes

The prerequisite for successful nuclear reprogramming involves the efficient uncoiling of the genomic DNA inside the nucleus that allows access to regulatory proteins. This is achieved through chromatin decondensation and histone modification that includes acetylation, methylation, and phosphorylation. In somatic cell nuclear transfer (SCNT), the histone modification pattern in the chromatin of a transplanted nucleus changes to that of the oocyte. These changes are enforced by transcription factors present in the oocyte cytoplasm. For example, oocyte-specific B4 or H1foo histone linkers replace the H1 histone and promote changes in gene expression patterns.

Yamanaka Factors

In 2006, Shinya Yamanaka discovered the Oct-4, Sox-2, Klf-4, and c-Myc transcription factors, referred to as Yamanaka factors, that are highly expressed in embryonic stem (ES) cells. The transcription factor transduction method reprograms the nucleus by overexpressing these transcription factors, thus inducing embryonic gene expression in somatic cells. These transformed cells are called induced pluripotent stem cells or iPSCs. The Oct-4 (octamer-binding transcription factor 4) is a protein encoded by the POU5F1 human gene. It plays a vital role in deciding the fate of the inner mass and embryonic stem cells and helps them to maintain pluripotency during embryonic development. The sex-determining region Y-box 2, referred to as Sox-2, plays an essential role in maintaining the self-renewal potential in ES cells. Klf-4 or Kruppel-like factor 4 belongs to the KLF family of zinc finger transcription factors and is majorly involved in the proliferation and differentiation of stem cells. c-Myc is the transcription factor belonging to the Myc family of protooncogenes, and it has an important role in cellular proliferation and metabolism.

 Core: Cell Biology

SI Units: 2019 Redefinition

JoVE 14504

Measurement is an indispensable part of analytical chemistry. The result of measurement helps quantify a substance's physical property and compare it with the physical property of another substance. Each measurement comprises two components - a number indicating the magnitude and a unit of measurement as a standard for comparison. Further, the same quantity can be measured using different units of measurement, which leads to differences in magnitude.

A standard set of units has been defined to help maintain consistency and avoid errors in scientific communication. These units are the 'Système International d'Unités' or SI units. The SI units for physical quantities such as mass, distance, temperature, time, electric current, luminous intensity, plane angle, solid angle, and amount of substance are termed fundamental units.

Other SI units are called derived SI units, constructed from the abovementioned fundamental SI units. For example, the SI unit for electric charge, Coulomb (C), is a product derived from multiplying the units of electric current (A) and time (s). Expression of quantities in non-SI units is also fairly common. For example, temperature is usually expressed in degrees Celsius (°C) or degrees Fahrenheit (°F) rather than the SI unit for temperature, Kelvin (K).

Quantities with very large or very small magnitude compared to a single SI unit are expressed in powers of ten. These powers of ten have prefixes that make their expression much easier, such as 'milli' for ten to the power of negative three or 'kilo' for ten to the power of positive three. This makes a kilometer the expression of a considerable distance, equivalent to a thousand meters or 103 meters.

 Core: Analytical Chemistry

Correlation and Regression

JoVE 14520

In statistics, correlation describes the degree of association between two variables. In the subfield of linear regression, correlation is mathematically expressed by the correlation coefficient, which describes the strength and direction of the relationship between two variables. The coefficient is symbolically represented by 'r' and ranges from -1 to +1. A positive value indicates a positive correlation where the two variables move in the same direction. A negative value suggests a negative correlation, where the two variables move in opposite directions. A value closer to +1 or −1 suggests that the two variables are strongly correlated, directly or inversely. A value close to zero implies no linear correlation between the two variables.

To obtain the correlation coefficient and the best-fit equation, statisticians use the method of linear regression. The best-fit equation can be used subsequently to predict the value of a signal ('y' in the linear equation) or calculate the concentration of the substance giving the signal ('x' in the linear equation).

 Core: Analytical Chemistry

Complexation Equilibria: Overview

JoVE 14536

Complexation reactions take place when dative or coordinate covalent bonds form between metal ions and ligands. The compounds formed in these reactions are called coordination compounds. The number of bonds formed between the metal ion and the ligands is called its coordination number. Generally, most metal ions in an aqueous solution are solvated by water molecules and thus exist as aqua complexes.

The equilibrium constant of the complexation reaction is represented as the formation constant Kf, also known as the stability constant Ks. Conversely, the dissociation of the complex is the reverse of the complex formation, and the dissociation equilibrium constant is called the dissociation constant Kd or instability constant Ki, which is equal to the reciprocal of the formation constant Kf

Complexation generally happens stepwise, where a metal ion initially complexes with one ligand, followed by the second ligand, until it satisfies its coordination number. The equilibrium constants for each reaction step are the stepwise formation constants represented as Kf1, Kf2, ...., Kfn. The overall formation or cumulative constant β is the product of the stepwise formation constants, or simply Kf.

The stepwise formation constant values decrease steadily due to statistical, coulombic, and steric factors. Statistical factors come into play after the first ligand attaches to the metal ion, which decreases the number of available sites for the next ligand and therefore decreases the probability of the next ligand binding. Coulombic factors also come into play when the metal ion complexes with the first ligand, because the positive charge on the metal ion decreases, thus decreasing the coulombic attraction for the next ligand. Steric factors have the most effect when a ligand is bulky, as the steric repulsions created by this ligand affect the ease with which the subsequent ligands interact with the metal ion.

 Core: Analytical Chemistry

EDTA: Direct, Back-, and Displacement Titration

JoVE 14575

The EDTA titration types for metal ion analysis include direct titration, back-titration, and replacement titration.

Direct titration involves buffering the metal ion solution to the desired pH and directly titrating with standard EDTA until the endpoint. The optimum pH ensures a large conditional formation constant of metal−EDTA and visibility of the free indicator color in the solution. In addition, auxiliary complexing reagents are used to prevent the precipitation of metal hydroxides and maintain the concentration of free metal ions in the solution.

Back titrations are useful when the metal ions block the indicators, react slowly, precipitate, or form inert complexes. This titration method adds a known excess of EDTA solution to the metal ion, and the solution is buffered to a desired pH. The excess EDTA is then back-titrated using a suitable indicator with a standard solution of a second metal ion to reach the endpoint.

Displacement or substitution titration is used when metal ions do not react adequately with the indicator. Since solochrome black is a poor indicator for directly titrating calcium ions, the calcium ions are first titrated with less stable excess [Mg(EDTA)]2− solution. The calcium ions displace all the magnesium ions to form a more stable [Ca(EDTA)]2− complex. The free magnesium ions released into the solution are equivalent to the calcium ions previously present. Finally, the free magnesium ion is titrated with a standard EDTA solution to obtain the endpoint.

 Core: Analytical Chemistry

Mass Analyzers: Common Types

JoVE 14591

The quadrupole mass analyzer consists of four cylindrical metal rods arranged in a diamond carrying a DC voltage and a radio-frequency AC voltage. The motion of ions through the quadrupole depends on the field strength, causing only ions of a certain m/z to resonate successfully and strike the detector at a given field strength. Though the transmission rate for these analyzers is high, the exact elemental composition of the sample is not determined because of low resolution; however, they are cost-effective. The typical resolution of a quadrupole analyzer is 1 Da (or 1 u).

In a time-of-flight (TOF) analyzer, the flight time of an ion from the source through the linear field-free drift tube to the detector is measured in microseconds, and this time is converted to an m/z value. The ions enter the drift tube with about the same amount of kinetic energy. Accordingly, ions with different m/z values travel with different velocities and reach the detector at different times. Lighter ions travel with higher velocity, resulting in a shorter time-of-flight and comparatively poor resolution. Resolution can also be impaired by differences in the amount of kinetic energy imparted to the ions, causing ions of the same m/z to have different velocities. Because of this, a reflectron (or ion mirror) is employed to increase the travel distance to the detector and slow down the ions. More energetic ions experience greater slowdown, so the reflectron corrects the kinetic energy spread and provides a better resolution overall. In addition, TOF analyzers offer high sensitivity.

Three-dimensional quadrupole ion-trap analyzers consist of a ring electrode and two endcap electrodes. A radio-frequency voltage traps the ions, which oscillate in the ring. Increasing the radio-frequency voltage stabilizes heavier ions and selectively ejects them to the detector via apertures in the endcaps. While other mass analyzers allow only a small fraction of ions to reach the detector, the three-dimensional quadrupole ion trap enables half of the ions to reach the detector, making these analyzers more sensitive. However, with a large number of ions stored in the ion trap, the space-charge effects result in decreased resolution.

 Core: Analytical Chemistry

Volatilization

JoVE 14624

Volatilization gravimetry is an analytical technique that measures the mass lost due to the volatilization of the substance. This technique is used to estimate the amount of volatile material in a sample. To perform this method, heat a known amount of the sample to a high temperature in a crucible or other suitable vessel. The volatile substance in the sample evaporates, and the vapor is completely expelled from the crucible either by heating the sample or bubbling a stream of inert gas through the vessel. The remaining non-volatile components are left in the crucible, and their mass is measured on a scale. The difference in mass before and after volatilization gives the mass of the volatile substance, which can then be used to calculate its concentration in the original sample. This approach helps determine the water of crystallization present in hydrated compounds.

Alternatively, the sample can be treated with a chemical reagent in a closed vessel to form a chemically different volatile species with a known composition. The volatilized substance is then selectively absorbed in a pre-weighed trap filled with a suitable absorbent. For example, the carbon dioxide released from carbonate-containing compounds can be selectively absorbed in soda-lime, and NaOH can be selectively absorbed in a trap with non-fibrous silicate. The traps can then be weighed to determine the amount of absorbed material.

 Core: Analytical Chemistry

Structure of Cardiac Muscles

JoVE 14853

Cardiac muscle, or myocardium, is a specialized type of muscle found exclusively in the heart. Its unique structural and functional characteristics enable the heart to perform its vital role of pumping blood throughout the body continuously and rhythmically. The cardiac muscle cells, or cardiomyocytes, possess an endomysium and perimysium but do not have an epimysium.

Compared to skeletal muscles, cardiac muscle cells are small and mostly have a single nucleus. Additionally, they are usually branched and composed of regularly arranged sarcomeres in the myofibrils, giving them a striated appearance. The T-tubules, which are short but broad, arise from the sarcolemma at the sarcomere boundaries and are closely associated with the sarcoplasmic reticulum. To sustain their high energy requirements, cells possess a multitude of mitochondria, as well as reserves of both glycogen and lipids. These energy stores are crucial for proper cellular function and enable cells to carry out their various metabolic processes with optimal efficiency.

Adjacent cardiac muscle cells are connected by intercalated discs, regions of thickened sarcolemma that interlock. These regions contain gap junctions that control the passage of ions and help propagate electrical signals during muscle contraction. The intercalated discs also have an abundance of desmosomes that anchor the cells to each other and prevent cell separation during contraction. This way, the cardiac cells can coordinate to produce the heart's rhythmic contractions.

Similar to skeletal muscle, cardiac muscle fibers can experience hypertrophy when there is an increase in workload. This results in a naturally enlarged heart, which is why many athletes have a larger heart. However, if the heart enlarges due to significant heart disease, it is considered pathological.

 Core: Anatomy and Physiology

Muscles of the Pelvic Floor and Perineum

JoVE 14874

The muscles of the pelvic floor and perineum are crucial for supporting the pelvic organs, controlling continence, and aiding in sexual function, childbirth, and core stability. They are typically divided into the superficial perineal layer and the deep pelvic floor layer.

Perineal Layer

The perineum is a diamond-shaped area below the pelvic diaphragm, divided into an anterior urogenital triangle that contains the external genitals and a posterior anal triangle housing the anus. The urogenital triangle consists of the superficial transverse perineal, bulbospongiosus, and ischiocavernosus muscles. The bulbospongiosus muscle surrounds the vaginal or penile areas to assist in sexual functions and support. The ischiocavernosus muscle helps maintain an erection in both males and females by compressing the outflow veins of the erectile tissues. Additionally, the superficial transverse perineal muscles support the pelvic organs and aid in maintaining continence.

Conversely, the posterior anal triangle contains the anus and is characterized by the presence of the external anal sphincter, which is essential for voluntary control over defecation. Surrounding this area are significant muscles like the obturator internus, which supports the pelvic organs and aids in lateral rotation of the thigh, and the inferior part of the levator ani from the pelvic diaphragm above, providing additional support and stability to the pelvic floor.

Pelvic Floor

Deeper within the pelvic floor lies a group of muscles that form a sling or hammock-like structure supporting the pelvic organs. The levator ani, a significant component of this layer, includes muscles such as the pubococcygeus, puborectalis, and iliococcygeus. These muscles are crucial for supporting the pelvic organs, maintaining continence, and facilitating childbirth by stretching and expanding as needed. The coccygeus muscle, extending from the ischial spine to the lower end of the sacrum and coccyx, complements the levator ani by providing additional support and stability.

 Core: Anatomy and Physiology

Graded Potential

JoVE 14891

Graded potentials are localized fluctuations in the cell membrane's electrical charge, commonly found in the dendrites of neurons. The magnitude of these potential changes depends on the strength of the initiating stimulus. In a membrane at its resting potential, a graded potential signifies a voltage shift either above -70 mV or below -70 mV.

Graded potentials fall into two categories: depolarizing and hyperpolarizing. Depolarizing graded potentials typically occur when sodium (Na+) or calcium (Ca2+) ions enter the cell. Since these ions have higher concentrations outside the cell and carry a positive charge, they flow inward, reducing the cell's negative charge relative to the extracellular environment.

Hyperpolarizing graded potentials result from either potassium (K+) exiting the cell or chloride (Cl-) entering it. When a positive charge exits the cell, the cell becomes more negatively charged. Conversely, the same effect occurs if a negative charge enters the cell.

Here's an explanation of graded potentials, their types, and their significance:

Types of Graded Potentials:

Excitatory Postsynaptic Potentials (EPSPs): These graded potentials depolarize the cell membrane, bringing it closer to the threshold for firing an action potential. EPSPs are typically caused by the influx of positively charged ions, such as sodium (Na+), through ligand-gated channels in response to neurotransmitter binding at synapses.

Inhibitory Postsynaptic Potentials (IPSPs): IPSPs are graded potentials that hyperpolarize the cell membrane, moving it away from the threshold for an action potential. IPSPs are often the result of the influx of negatively charged ions, such as chloride (Cl-), or the efflux of potassium (K+).

Significance of Graded Potentials:

Information Processing: Graded potentials are essential for information processing in the nervous system. They occur at synapses, where neurons communicate with each other. EPSPs and IPSPs integrate information from multiple inputs, allowing the neuron to decide whether to generate an action potential and transmit a signal.

Spatial Summation: Neurons can receive inputs from many synapses. Graded potentials from different synapses can sum together spatially, either reinforcing or inhibiting each other. This spatial summation helps determine whether the neuron reaches the threshold for firing an action potential.

Temporal Summation: Graded potentials can also sum over time, which is called temporal summation. When multiple graded potentials occur in rapid succession, their effects can add up, potentially leading to the generation of an action potential.

In summary, graded potentials are essential for fine-tuning information processing in the nervous system. They allow neurons to integrate excitatory and inhibitory signals from various sources, making decisions about whether to transmit signals further down the neural circuit. This ability to modulate and integrate signals is critical for the complexity and functionality of the nervous system.

 Core: Anatomy and Physiology

Motor and Sensory Areas of the Cortex

JoVE 14907

The cerebral cortex, the brain's outermost layer, is pivotal in processing complex cognitive tasks, emotions, and various sensory inputs and executing voluntary motor activities. This intricate structure is divided into three primary functional areas: the motor areas, sensory areas, and association areas.

Motor Areas

The motor areas located in the frontal lobe are central to controlling voluntary movements. This region is further subdivided into the primary motor cortex and the premotor cortex. The primary motor cortex, situated in the precentral gyrus, is directly involved in controlling skeletal muscle movements. It has a somatotopic organization, meaning that different body parts are represented in specific areas of the cortex, known as the motor homunculus. Neurons within this cortex send projections to specific muscles, orchestrating precise movements ranging from the delicate motions of the fingers to the complex coordination required for speech or walking.

The premotor cortex is adjacent to the primary motor cortex and anterior to the precentral gyrus. This area is involved in planning and coordinating movements, working closely with the primary motor cortex to execute complex motor tasks. A specialized part of the premotor cortex, Broca's area, located near the lateral sulcus, plays a pivotal role in speech production. It controls the muscles involved in speech, demonstrating the specialized functions certain cortical areas possess.

Sensory Areas

The sensory areas of the cerebral cortex are intricately designed to process and interpret sensory information from both external and internal environments. Each area specializes in handling different sensory inputs — touch, sight, hearing, smell, and taste. This specialization facilitates a rich and nuanced perception of the world.

The primary somatosensory cortex, situated in the postcentral gyrus of the parietal lobe, plays a fundamental role in processing tactile information from the body. Like the primary motor cortex, this area is also organized somatotopically. This organization allows for precise spatial discrimination, enabling individuals to identify where a particular touch, pressure, or pain stimulus is coming from.

The primary visual cortex is located in the occipital lobe and is the main gateway for visual processing. It receives input from the retina of the eye through the optic nerves and pathways. The primary visual cortex is responsible for decoding basic visual information, such as light intensity, color, and movement, before this information is sent to other visual association areas for further processing and interpretation.

The primary auditory cortex, situated in the temporal lobe, is crucial for processing sounds. It receives input from the cochlea of the inner ear through the auditory pathways. This cortex is organized tonotopically, meaning different areas process different sound frequencies.

The olfactory cortex, located in the temporal lobe, is involved in the sense of smell. It receives information directly from the olfactory bulbs, which pick up chemical signals from the nose. The olfactory cortex plays a role in identifying different odors, and it is closely linked to memory and emotions, explaining why certain smells can evoke strong memories or feelings.

Finally, the primary gustatory cortex, found in the insula, is responsible for taste perception. It processes information from the taste buds on the tongue, differentiating between sweet, salty, sour, bitter, and umami (savory) tastes.

Association Areas

Beyond the motor and sensory areas lie the association areas, which integrate the sensory and motor information to support complex cognitive functions such as learning, memory, reasoning, and emotions. These areas are not dedicated to processing simple sensory stimuli or motor commands; instead, they are involved in the interpretation, planning, and execution of tasks that require a higher level of cognitive function. For example, they play a role in recognizing objects and faces, understanding language, and planning future actions.

 Core: Anatomy and Physiology

Cranial Nerves: Types Part II

JoVE 14923

Cranial nerves are responsible for transmitting motor and sensory information between the brain and various parts of the body. There are twelve pairs of cranial nerves. While the first six innervate the head and neck, the latter six nerves innervate the head and neck, as well as organs and tissues in the thoracic and abdominal cavities. They facilitate communication, expression, and autonomic control within the human body.

Facial Nerve (Cranial Nerve VII)

Cranial nerve VII, or the facial nerve, is a mixed nerve responsible for facial expressions, taste sensation, and salaivary gland control. It emerges from the brainstem and travels through the temporal bone before branching out to innervate the muscles of facial expression. Additionally, it carries taste sensations from the anterior two-thirds of the tongue and controls secretions from the lacrimal glands and salivary glands.

Vestibulocochlear Nerve (Cranial Nerve VIII)

The vestibulocochlear nerve, or cranial nerve VIII, is a sensory nerve dedicated to hearing and balance. It consists of two components — the cochlear nerve, which transmits sound information from the inner ear to the brain, and the vestibular nerve, which conveys information about balance and spatial orientation. Proper functioning of the vestibulocochlear nerve is crucial for auditory perception and maintaining equilibrium.

Glossopharyngeal Nerve (Cranial Nerve IX)

The glossopharyngeal nerve has both motor and sensory functions. It innervates the pharynx, contributing to swallowing, and carries taste and sensory information from the posterior one-third of the tongue. The glossopharyngeal nerve also plays a role in monitoring oxygen and carbon dioxide levels in the blood, illustrating its involvement in autonomic control.

Vagus Nerve (Cranial Nerve X)

The vagus nerve, or cranial nerve X, is the longest cranial nerve and remarkably extends beyond the head and neck to innervate the thorax and abdomen. It is a mixed nerve that plays a vital role in the autonomic nervous system by controlling the heart, lungs, and digestive tract. The functions of the vagus nerve include speech, coughing, and gastrointestinal motility.

Accessory Nerve (Cranial Nerve XI)

The accessory nerve is the eleventh cranial nerve. It is primarily a motor nerve that innervates the sternocleidomastoid and trapezius muscles. These muscles are integral to head movement and shoulder elevation, enabling movements essential for various physical activities, from nodding to lifting objects.

Hypoglossal Nerve (Cranial Nerve XII)

The twelfth cranial nerve, known as the hypoglossal nerve, is solely a motor nerve that governs the movements of the tongue. It is crucial for speech, chewing, and swallowing. The hypoglossal nerve's control over the tongue's movements is indispensable for articulation and the physical eating process, impacting communication and nutrition.

 Core: Anatomy and Physiology

Brain Waves

JoVE 14942

Brain waves are electrical signals generated by the neurons in the brain, which are regularly monitored to measure mental activities. Brain waves and their frequency ranges can be measured using an electroencephalogram or EEG. There are four main types of brain waves, each with distinct characteristics:

  • • Alpha Waves – Alpha brain waves, between 8 and 13 Hz ( Hertz), occur when a person is relaxed with eyes closed. The alpha waves then disappear when one is concentrating on specific tasks.
     
  • • Beta Waves – Beta brain waves, between 14 and  30 Hz, occur when an individual is awake and alert. This type of brainwave is typical when engaged in problem-solving or planning.
  • • Theta Waves – Theta brain waves, from 4 to 7 Hz, occur primarily in children. EEG of adults occasionally shows theta waves during periods of emotional stress. Their presence in adults not under any stress indicates the presence of tumors or other pathological conditions affecting the brain.
  • • Delta Waves – Delta brain waves have a shallow frequency at 1 to 5 Hz and are seen during an adult's deepest level of sleep. They are also seen in infants. Delta waves in wakeful adults may also indicate brain damage.

An EEG can detect various types of disorders, including epilepsy, dementia, stroke, and traumatic brain injury (TBI). It can also be used to monitor for changes in an individual's mental state over time. The test is usually performed by placing electrodes on the scalp, which detect electrical signals from the brain and transmit them to a computer for analysis.

 Core: Anatomy and Physiology

Accessory Structures of the Eye

JoVE 14961

Optical perception, or vision, is an extraordinary sense dependent on converting light signals received via the ocular organs. These organs, known as eyes, are securely positioned within the bony cavities of the skull, called orbits. The orbits serve a dual purpose: a protective shield for the ocular globes and a stable attachment point for the soft ocular tissues. The eye's external protective mechanisms include the eyelids, which are edged with lashes that act as a barrier against foreign particles potentially causing abrasions. Covering the inner surface of each eyelid is a delicate membrane termed the palpebral conjunctiva. This membrane covers the sclera, the white portion of the eye, forming a connection between the eyelids and the eyes. The lacrimal gland, located within the orbit at a position superior and lateral to the eye, is responsible for the production of tears. These tears traverse through the lacrimal duct to the eye's medial corner, moistening the conjunctiva and effectively cleansing the eye of any foreign bodies.

The ocular movements within the orbital cavity are engineered by contracting six periorbital muscles originating from the periorbital bone structure and attaching to the spherical ocular structure. These muscles comprise the superior, medial, inferior, and lateral rectus, strategically positioned at the major directional points around the ocular organ. Upon contraction, the respective muscle guides the eye toward the direction of the contracting muscle. For instance, an upward gaze is facilitated by the contraction of the superior rectus. The superior oblique muscle, attached obliquely to the superior surface of the eyeball, is located posteriorly in the orbital cavity near the origination point of the rectus quartet. Contraction of the superior oblique muscles results in a lateral eyeball rotation. The inferior oblique muscle, originating from the orbital floor and attaching to the inferior-lateral surface of the eye, executes a similar lateral rotation of the eye upon contraction, acting in opposition to the superior oblique. This ocular rotation facilitated by the oblique duo is requisite due to the slight eye misalignment on the sagittal plane. Thus, in situations requiring upward or downward gaze, a fine rotational adjustment is necessary to counteract the off-center pull of the superior and inferior rectus muscles. The contraction of the inferior oblique muscle manages this adjustment. Another muscle within the orbital architecture, the levator palpebrae superioris, is tasked with elevating and retracting the upper eyelid. This movement usually synchronizes with the elevation of the eye by the superior rectus. Three cranial nerves handle the innervation of these extraocular muscles. Specifically, the lateral rectus (which executes the abduction of the eye) is innervated by the abducens nerve; the superior oblique receives signals from the trochlear nerve; and the remainder of the muscles, including the levator palpebrae superioris, are innervated by the oculomotor nerve. The motor nuclei of these cranial nerves maintain connections with the brain stem, which, in turn, coordinates the ocular movements.

Diseases of accessory structures of the eye: 

Diseases of accessory structures of the eye include specific allergic reactions, bacterial conjunctivitis, blepharitis, and dry eye syndrome:

  • • Allergic reactions involve increased histamine production by the body's immune cells, resulting in redness, itching, and swelling around the eyes.
  • • Bacterial conjunctivitis is caused by bacteria lodged within the palpebral conjunctiva and is characterized by redness, swelling, and discharge in the eye.
  • • Blepharitis occurs due to inflammation of multiple meibomian glands involving both eyelids. It can cause itching, a burning sensation in the eyes, scaling of the eyelid margins, and redness.
  • • Dry eye syndrome is caused by the decreased production or excessive evaporation of tears, leading to dryness, burning, irritation, and eye redness.

Treatment for these conditions include

  • • Artificial tears for dry eye syndrome
  • • Corticosteroids and antibiotics for bacterial conjunctivitis
  • • Antihistamines for allergies
  • • In cases of blepharitis, eyelid hygiene, which involves gently applying warm compresses, massaging the eyelids, and using an eyelid scrub, is recommended. These treatments can help to reduce inflammation and prevent the progression of the disease.

The eyes also play a significant role in facial expressions. Controlling ocular movements is essential for expressing surprise, joy, or displeasure. For example, when surprised, our eyes tend to open wider, whereas an upward and outward movement of the eyes characterizes a smile. The ocular muscles, with their fine-tuned control of eye movements and adjustments in a degree of aperture, are instrumental in conveying emotions to observers. This means that proper functioning of these muscles is essential for regular facial expressions. Disease conditions such as strabismus can impair this ability, causing a loss of emotive facial control. Therefore, regular eye exams are essential to ensure the health of the ocular organs and their proper functioning for both visual acuity and emotional expression. Additionally, wearing protective eyewear while engaging in outdoor activities is recommended to guard against potential injuries to the delicate tissues surrounding the eyes. By doing so, we can ensure our ocular organs' long-term health and associated functions.

In addition, it is crucial to be aware of eye strain. Eye strain is caused by excessive or prolonged use of the eyes for activities such as computer work, reading, or writing. Symptoms of eye strain include headaches, blurred vision, dry eyes, neck and shoulder aches, and tiredness. To prevent eye strain, taking regular breaks from tasks requiring visual concentration, such as taking a break every 30 minutes, is essential. During these breaks, it is recommended that you look away from the screen or book at a distant point to rest your eyes. It is also advisable to use proper lighting and adjust the contrast and brightness levels on screens to reduce strain. Lastly, it is recommended to use eye drops or lubricant ointment if needed to reduce dryness and irritation in the eyes. By following these simple preventative measures, we can help ensure the long-term health of our eyes and vision.

 Core: Anatomy and Physiology

Target Cell Response to Hormones

JoVE 14977

Hormones intricately bind to receptors on the surface or within target cells, initiating a cascade of cellular responses.

Notably, the cellular response can be regulated by altering the number of receptors expressed in the cell. For example, prolonged exposure to elevated hormone levels results in a gradual decline or down-regulation in the number of receptors for that specific hormone on the cell surface. Conversely, in response to low hormone levels, cells may use up-regulation, producing an increased quantity of a particular receptor to enhance cellular sensitivity.

The interplay between hormones can lead to diverse cellular responses. Certain hormones exhibit permissive interactions, where the presence of one hormone enables another to exert its effects synergistically. An example is the collaboration between epinephrine and thyroid hormones, where thyroid hormones facilitate the effective stimulation of lipolysis by epinephrine in target cells. Furthermore, hormones like follicle-stimulating hormone and testosterone demonstrate synergistic action, working in concert to elicit an amplified response, such as the normal production of sperm. On the contrary, some hormones engage in antagonistic relationships, resulting in opposing cellular responses. A classic illustration is the interaction between insulin and glucagon, where insulin stimulates a decrease in blood glucose levels while glucagon acts to increase them. This intricate network of interactions underscores the dynamic nature of hormonal regulation and its profound influence on cellular physiology.

 Core: Anatomy and Physiology

Diabetes Mellitus: Overview and Type I Subtype

JoVE 14993

Diabetes mellitus is a chronic metabolic disorder characterized by high blood glucose levels due to inadequate insulin production, insulin resistance, or both. The condition affects millions worldwide and can significantly impact their health and quality of life.

Type 1 diabetes is an autoimmune disease in which the immune system mistakenly attacks and destroys the insulin-producing beta cells in the pancreas. As a result, the body is unable to produce sufficient insulin, and individuals with type 1 diabetes require lifelong insulin replacement therapy.

On the other hand, type 2 diabetes is primarily caused by insulin resistance, a condition in which the body's cells become less responsive to insulin. As a result, the reduced ability of insulin to facilitate glucose uptake leads to elevated blood glucose levels. Type 2 diabetes is often associated with lifestyle factors such as obesity, physical inactivity, and poor dietary choices. It is a progressive condition that may initially be managed with lifestyle modifications and oral medications, but some individuals may eventually require insulin therapy.

Uncontrolled diabetes can lead to a range of complications affecting various organs and systems in the body. These complications include cardiovascular disease, kidney disease, nerve damage or neuropathy, retinopathy, and impaired wound healing. It is important to manage blood glucose levels through medication, dietary changes, regular physical activity, and regular monitoring to prevent or delay the onset of complications.

Diabetes treatment aims to maintain blood glucose levels within the prescribed range. Lifestyle modifications like a healthy diet, monitoring blood glucose levels, and exercise help effective diabetes management. Additionally, oral medications to improve insulin sensitivity or stimulate insulin production and insulin injections or pumps for individuals with type 1 or advanced type 2 diabetes are prescribed.

 Core: Anatomy and Physiology

Maximum Power Transfer

JoVE 15055

Numerous practical applications within engineering disciplines, such as telecommunications, necessitate optimizing power delivery to a connected load. This pursuit, however, entails inherent internal losses, which can either equal or exceed the power supplied to the load. The Thevenin equivalent circuit is helpful in finding the maximum power a linear circuit can deliver to a load. It is assumed in this context that the load resistance can be adjusted.

By substituting the entire circuit with its Thevenin equivalent while preserving the load as shown in Figure 1,

Figure1

Figure 1:Circuit used for maximum power transfer

 the power delivered to the load is given by

Equation1.................(1)

For a given circuit, Thevenin equivalent resistance and voltage are fixed. The variation of power delivered to the load with load resistance is small for small or large values of load resistance but maximum for some value load resistance between zero and infinity. The maximum power occurs when the load resistance equals the Thevenin's resistance expressed as

Equation2.................(2)

This is known as the maximum power transfer theorem. Using equations 1 and 2, the expression for maximum power transferred is obtained and is expressed as

Equation3..................(3)

This Equation applies only when equation 2 is valid, otherwise, the power delivered to the load is determined using Eq. 1.

 Core: Electrical Engineering

Inverting and Non-inverting OpAmps

JoVE 15078

In an inverting amplifier, the input voltage is connected through a resistor to the inverting terminal. Meanwhile, the non-inverting terminal is grounded and a feedback resistor is established between the inverting and output terminal, as depicted in Figure 1.

Figure1

Figure 1: The inverting amplifier

The objective is to discern the relationship between the input voltage (vi) and the output voltage (vo). By applying Kirchhoff's Current Law (KCL) and presuming the operational amplifier to be ideal, the expression for gain is obtained. An inverting amplifier has the ability to reverse the polarity of the input signal while simultaneously amplifying it. It is worth noting that the gain depends solely on the external elements connected to the op-amp, with the gain being the feedback resistance divided by the input resistance.

Another critical application of the op-amp is seen in the noninverting amplifier, as illustrated in Figure 2.

Figure2

Figure 2: The non-inverting amplifier

Here, the input voltage (vi) is applied directly at the non-inverting input terminal, while resistor R is connected between the ground and the inverting terminal. The focus here is on the output voltage and the voltage gain. Applying KCL at the inverting terminal and substituting the voltage values gives an expression for voltage gain. If the feedback resistor gets short-circuited or if the input resistor is open-circuited, the gain transforms into unity, forming a voltage follower or unity gain amplifier. Such a circuit possesses a high input impedance, making it useful as an intermediate-stage (or buffer) amplifier to segregate one circuit from another.

 Core: Electrical Engineering

RL Circuit without Source

JoVE 15095

When a DC source is suddenly disconnected from an RL (Resistor-Inductor) circuit, the circuit becomes source-free. Assuming the inductor has an initial current denoted as I0, the initial energy stored in the inductor can be determined.

Applying Kirchhoff's voltage law around the loop of the circuit and substituting the voltages across the inductor and resistor yields a first-order differential equation. A logarithmic equation is obtained by rearranging the terms in this equation, integrating it, and applying the limits. Taking the exponential on both sides of this equation yields the final expression of the circuit's natural response.

If the current is plotted versus time, an exponential decrease is observed in the initial current. This behavior can be expressed in terms of the time constant, which for an RL circuit is the ratio of inductance to resistance. This time constant represents the speed at which the circuit responds to changes in the input signal.

By using the current expression, the voltage across the resistor and the power that gets dissipated in the resistor can be calculated. The power dissipated is essentially the rate at which energy is lost in the form of heat.

To find the total energy absorbed by the resistor, the power dissipated is integrated over time. As time approaches infinity, the energy absorbed by the resistor approaches the initial energy stored in the inductor. This implies that the initial energy stored in the inductor gradually dissipates in the resistor until the inductor's energy is depleted.

In conclusion, understanding the behavior of RL circuits when the DC source is removed offers valuable insights into the transient response of these circuits. This knowledge is fundamental for designing and analyzing circuits in applications such as power electronics and communication systems, where inductors are used extensively to filter or shape signals.

 Core: Electrical Engineering

Source Transformation for AC Circuits

JoVE 15112

The process of source transformation in the frequency domain entails the conversion of a voltage source, positioned in series with an impedance, into a current source that is parallel to an impedance, or the other way around. It is essential to maintain the following relationships while transitioning from one source type to another.

Equation1

Equation2

In order to determine the unknown voltage for a circuit composed of a current source and a collection of resistors, capacitors, and inductors - each with their distinct known impedance, a series of steps are followed. Initially, the voltage source is converted into a current source, and the values of the source current (Is) and impedance (Zs) are established.

Subsequently, transforming the current source back to a voltage source results in a different circuit. From this derived circuit, the source voltage (Vs) is calculated using the previously determined values. Finally, by applying the voltage division rule, the unknown voltage across the resistance can be identified.

 Core: Electrical Engineering

Titration of a Strong Acid with a Strong Base

JoVE 17360

During the titration of a strong acid with a strong base, pH calculations are primarily based on the concentration of residual hydronium or hydroxide ions. Initially, a strong acid like hydrochloric acid fully dissociates, creating hydronium and chloride ions, resulting in a low pH. The addition of a strong base like sodium hydroxide alters the concentration of hydronium ions by neutralizing them. As more base is added, the pH gradually increases. At the equivalence point, all hydronium ions are neutralized, achieving a neutral pH. Beyond this point, the excess hydroxide ions dictate the pH, making the solution basic. The new pH can be calculated by considering the concentration of these excess hydroxide ions.

 Core: Analytical Chemistry

Drosophila melanogaster Embryo and Larva Harvesting and Preparation

JoVE 5094

Drosophila melanogaster embryos and larvae are easy to manipulate and develop rapidly by mechanisms that are analogous to other organisms, including mammals. For these reasons, many researchers utilize fly embryos and larvae to answer questions in diverse fields ranging from behavioral to developmental biology. Prior to experimentation, however, the embryos and larvae must first be collected.

This video will first demonstrate how "egg-laying cups" are used to collect Drosophila embryos on agar plates. The harvest and dechorionation of embryos will then be described. Next, the video will demonstrate how to identify and manipulate Drosophila in one of the three larval stages that follow the embryo stage. Finally, examples of some of the ways in which fly embryos and larvae are used in biological research are provided.

 Biology I

In ovo Electroporation of Chicken Embryos

JoVE 5156

Electroporation is a technique used in biomedical research that allows for the manipulation of gene expression via the delivery of foreign genetic material into cells. More specifically, in ovo electroporation is performed on early developing chicks (Gallus gallus domesticus) contained within their eggshells. In this procedure, DNA or knockdown constructs are first injected into a target tissue. However, the genetic material is unable to penetrate the plasma membrane to carry out its function within the cell. To solve this problem, an electrical field is applied, causing temporary disruptions to membrane stability. This electric field also causes the negatively charged nucleic acids to migrate toward the positively charged electrode through the holes in the plasma membrane, thus effectively driving the DNA or knockdown construct into the cell. The major advantage of this technique is that the delivery of genetic material can be localized to isolated cell types at specific developmental time points. As a result, the genetic mechanisms that govern individual developmental events can be examined.

This video provides an overview of the principles behind in ovo electroporation and introduces the tools required for the technique, including capillary needles, electrodes, and an electroporator. A step-by-step protocol for carrying out the procedure is also presented prior to discussion of a few fascinating examples of how the technique is used to perform a variety of genetic manipulations in chicken embryos.

 Biology II

Explant Culture of Neural Tissue

JoVE 5209

The intricate structure of the vertebrate nervous system arises from a complex series of events involving cell differentiation, cell migration, and changes in cell morphology. Studying these processes is essential to our understanding of nervous system function as well as our ability to diagnose and treat disorders that result from abnormal development. However, neural tissues are relatively inaccessible for experimental manipulations, especially in embryonic mammals. As a result, many scientists take advantage of explant culture in order to study neurodevelopmental processes in an “organotypic” environment, meaning that the tissue is removed from the organism but its complex cellular architecture is maintained. Generally, explant cultures are created by careful dissection of neural tissue that is then submerged in carefully designed growth media and cultured in vitro.

This video will first provide a brief overview of neural explant culture, including its advantages over other in vitro methods and important considerations for maintaining healthy tissue. Next, a general protocol will be provided for setting up an explant culture from embryonic mouse brain, outlining the isolation of embryos from the mother and dissection of the brain. The presentation also includes an overview of slice culture, in which thin sections of nervous system tissue are generated for improved visual access to the developing cells. Lastly, a few applications of these techniques will be provided to demonstrate how they can be used to answer important questions in the neurodevelopmental field.

 Neuroscience

An Introduction to Stem Cell Biology

JoVE 5331

Cells that can differentiate into a variety of cell types, known as stem cells, are at the center of one of the most exciting fields of science today. Stem cell biologists are working to understand the basic mechanisms that regulate how these cells function. These researchers are also interested in harnessing the remarkable potential of stem cells to treat human diseases.

Here, JoVE presents an introduction to the captivating world of stem cell biology. We begin with a timeline of landmark studies, from the first experimental evidence for hematopoietic stem cells in the 1960s, to more recent breakthroughs like induced pluripotent stem cells. Next, key questions about stem cell biology are introduced, for example: How do these cells maintain their unique ability to undergo self-renewal? This is followed by a discussion of some prominent methods used to answer these questions. Finally, several experiments are presented to demonstrate the use of stem cells in regenerative medicine.

 Developmental Biology

Eye Tracking in Cognitive Experiments

JoVE 5421

Eye tracking as the name suggests involves tracking of eye-movements. It is a non-invasive, sensitive tool that quantifies and measures eye-movements to describe an individuals' cognitive state. An eye-movement between two fixation points is called a saccade, which is one of the fastest motor movements in our body. By observing the profiles of these eye movements, scientists can better understand neural deficits in patients with cognitive impairments.

In this video, we will first look at an overview of different eye movements that eye tracking can capture and the type of data that can be collected. Then, the basic setup and experimental design are reviewed, including different types of eye trackers and details to optimize the eye tracking equipment. Finally, we will take a look at a few specific experiments utilizing eye tracking as a tool to study cognition.

 Behavioral Science

Assembly of a Reflux System for Heated Chemical Reactions

JoVE 5516

Source: Laboratory of Dr. Philip Miller — Imperial College London

Many chemical experiments require elevated temperatures before any reaction is observed, however heating solutions of reactants can lead to loss of reactants and/or solvent via evaporation if their boiling points are sufficiently low. In order to ensure no loss of reactants or solvent, a reflux system is used in order to condense any vapors produced on heating and return these condensates to the reaction vessel. 

 Organic Chemistry

Recombineering and Gene Targeting

JoVE 5553

One of the most widely used tools in modern biology is molecular cloning with restriction enzymes, which create compatible ends between DNA fragments that allow them to be joined together. However, this technique has certain restrictions that limit its applicability for large or complex DNA construct generation. A newer technique that addresses some of these shortcomings is recombineering, which modifies DNA using homologous recombination (HR), the exchange between different DNA molecules based on stretches of similar or identical sequences. Together with gene targeting, which takes advantage of endogenous HR to alter an organism’s genome at a specific loci, HR-based cloning techniques have greatly improved the speed and efficacy of high-throughput genetic engineering.

In this video, we introduce the principles of HR, as well as the basic components required to perform a recombineering experiment, including recombination-competent organisms and genomic libraries such as bacterial artificial chromosomes (BAC). We then walk through a protocol that uses recombineering to generate a gene-targeting vector that can ultimately be transfected into embryonic stem cells to generate a transgenic animal. Finally, several applications that highlight the utility and variety of recombineering techniques will be presented.

 Genetics

FM Dyes in Vesicle Recycling

JoVE 5648

FM dyes are a class of fluorescent molecules that has found important use in studying the vesicle recycling process. By virtue of a chemical structure, these molecules can insert themselves into the outer leaflet of phospholipid bilayer membranes. After membrane insertion, they are internalized into the cell via endocytosed vesicles, and released when these vesicles recycle back to the membrane. Since, these dyes fluoresce strongly in the hydrophobic environment within membranes and weakly in the extracellular compartment, FM fluorescence levels can be used to track vesicular activity throughout the recycling process.

This video provides an introduction to the use of FM dyes in experiments aimed to examine vesicle recycling. We first review the biochemistry of FM dyes and how their properties permit their use in these experiments. We then go through a general protocol for using FM dyes in such studies, and finally, discuss some recent research that makes use of these unique molecules.

 Cell Biology

Tandem Mass Spectrometry

JoVE 5690

In tandem mass spectrometry a biomolecule of interest is isolated from a biological sample, and then fragmented into multiple subunits in order to help elucidate its composition and sequence. This is accomplished by having mass spectrometers in series. The first spectrometer ionizes a sample and filter ions of a specific mass to charge ratio. Filtered ions are then fragmented and passed to a second mass spectrometer where the fragments are analyzed.

This video introduces the principles of tandem mass spectrometry, including mass-to-ratio selection and dissociation methods. Also shown is a general procedure for analyzing a biochemical compound using tandem mass spectrometry with collision-induced dissociation. The applications section covers selection reaction monitoring, determination of protein post-translation modifications, and detection of tacrolimus levels in blood.

Tandem mass spectrometry links together multiple stages of mass spectrometry to first isolate a biomolecule, and then determine aspects of its chemical makeup. Biomolecules have large, complex structures, making it difficult to determine their molecular composition. Tandem mass spectrometry selects a molecule of interest that is later fragmented into multiple subunits, which can help elucidate its identification and sequence. This video will show the concepts of tandem mass spectrometry, a general procedure, and some of its uses in biochemistry.

Tandem mass spectrometry begins as a typical mass spec instrument: with an ion source, which converts the sample into ions, and a mass analyzer, which separates the ions based on their mass-to-charge ratio. A common mass analyzer, the quadrupole, only allows ions with a specific ratio through, while the others crash into the rods of the apparatus. The species allowed through, called the precursor ion, is the biomolecule of interest. The ion moves into a collision cell, typically another quadrupole, where energy is applied to fragment the ion in a predictable pattern.

These fragments move into another mass analyzer, such as a time-of-flight, which separates these "product ions". The product ions are then sent to the detector, as in a normal MS instrument. In the case of an unknown protein, the resulting spectrum contains numerous overlapping fragments, making a definitive complete sequence of the biomolecule difficult to generate. However, the spectral pattern is unique for a given protein. Analysis software compares the spectrum to a database of known peptide sequences, elucidating the unknown protein from the overlapping fragments.

Depending on the sample and desired degree of fragmentation, multiple fragmentation methods are possible. Fragmentation patterns depend on how the energy is transferred, its amount, and how it is distributed through the precursor ion. Energy can be transferred via neutral particles, radiation, or electrons. Using neutral atoms, a process called collision-induced dissociation or CID, primarily cleaves at the peptide bond between the amino acids, ideal for their identification.

Now that the basics of the technique have been covered, let's look at CID tandem mass spectrometry being used to study a component of bacterial cell envelopes.

As with all mass spectrometric experiments, the first step is to ionize the sample. For biomolecules, this is typically done with matrix assisted laser desorption or electrospray ionization. The precursor ion signal is then optimized by tuning of the ion optics. Once done, the target is isolated and the fragmentation method is chosen, such as CID.

The strength of an applied voltage, which accelerates the precursor ion into the collision cell, affects the degree of fragmentation. This voltage is increased until the precursor is roughly 10% abundance compared to the highest product ion. Multiple spectra are acquired and averaged until a sufficient signal-to-noise ratio is achieved. The number of scans needed is dependent on the signal intensity of the original precursor ion and can range from 3 to 300.

The analyte in this example, lipid A from Escherichia coli K-12, had 19 major fragments after CID. Lipid A's general structure is well known, allowing software to reconstruct the specific composition from the sample.

Now that we've looked the procedure, let's look at some of the ways tandem mass spectrometry is used in biochemistry.

A common scanning mode in tandem mass spectrometry is selected reaction monitoring, or SRM. In SRM, both mass analyzers are fixed to a selected mass-to-charge ratio, focusing on specific precursor and product ions. Because of SRM's high degree of sensitivity, the spectra of peptide standards of known concentration can be utilized and compared to that of the unknown samples, allowing proteins of interest to be quantified.

Proteins are commonly modified after translation, typically by the addition of functional groups such as methyl groups, phosphate groups, or sugars, known as glycans. These are important in cell signaling processes, elucidating how cells communicate with one another. Because tandem mass spectrometry fragments the proteins into smaller components, it is possible to determine the location of the PTM to the specific fragment or even an amino acid. Some modifications, such as acetylation and trimethylation, are difficult to differentiate by mass alone, so chromatographic separation is performed before the mass spectrometry.

Many analytes in patient's blood are found at concentrations below the limit of detection for typical mass spectrometry. Another advantage of SRM is that it discards all but one product ion, increasing the sensitivity and enhancing the lower detection limit by up to 100 fold. In this example, the immunosuppressant drug, tacrolimus, could be detected at levels of 1 ng/mL.

You've just watched JoVE's video on tandem mass spectrometry. This video described the theory of the instrument, went over a general procedure, and explained some of the ways the technique is currently being utilized. Thanks for watching!

 Biochemistry